Text-to-image¶
The Stable Diffusion model was created by researchers and engineers from CompVis, Stability AI, Runway, and LAION. The StableDiffusionPipeline
is capable of generating photorealistic images given any text input. It's trained on 512x512 images from a subset of the LAION-5B dataset. This model uses a frozen CLIP ViT-L/14 text encoder to condition the model on text prompts. With its 860M UNet and 123M text encoder, the model is relatively lightweight. Latent diffusion is the research on top of which Stable Diffusion was built. It was proposed in High-Resolution Image Synthesis with Latent Diffusion Models by Robin Rombach, Andreas Blattmann, Dominik Lorenz, Patrick Esser, Björn Ommer.
The abstract from the paper is:
By decomposing the image formation process into a sequential application of denoising autoencoders, diffusion models (DMs) achieve state-of-the-art synthesis results on image data and beyond. Additionally, their formulation allows for a guiding mechanism to control the image generation process without retraining. However, since these models typically operate directly in pixel space, optimization of powerful DMs often consumes hundreds of GPU days and inference is expensive due to sequential evaluations. To enable DM training on limited computational resources while retaining their quality and flexibility, we apply them in the latent space of powerful pretrained autoencoders. In contrast to previous work, training diffusion models on such a representation allows for the first time to reach a near-optimal point between complexity reduction and detail preservation, greatly boosting visual fidelity. By introducing cross-attention layers into the model architecture, we turn diffusion models into powerful and flexible generators for general conditioning inputs such as text or bounding boxes and high-resolution synthesis becomes possible in a convolutional manner. Our latent diffusion models (LDMs) achieve a new state of the art for image inpainting and highly competitive performance on various tasks, including unconditional image generation, semantic scene synthesis, and super-resolution, while significantly reducing computational requirements compared to pixel-based DMs. Code is available at https://github.com/CompVis/latent-diffusion.
Tip
Make sure to check out the Stable Diffusion Tips section to learn how to explore the tradeoff between scheduler speed and quality, and how to reuse pipeline components efficiently!
If you're interested in using one of the official checkpoints for a task, explore the CompVis, Runway, and Stability AI Hub organizations!
mindone.diffusers.StableDiffusionPipeline
¶
Bases: DiffusionPipeline
, StableDiffusionMixin
, IPAdapterMixin
, TextualInversionLoaderMixin
, LoraLoaderMixin
, FromSingleFileMixin
Pipeline for text-to-image generation using Stable Diffusion.
This model inherits from [DiffusionPipeline
]. Check the superclass documentation for the generic methods
implemented for all pipelines (downloading, saving, running on a particular device, etc.).
The pipeline also inherits the following loading methods
- [
~loaders.TextualInversionLoaderMixin.load_textual_inversion
] for loading textual inversion embeddings - [
~loaders.LoraLoaderMixin.load_lora_weights
] for loading LoRA weights - [
~loaders.LoraLoaderMixin.save_lora_weights
] for saving LoRA weights - [
~loaders.FromSingleFileMixin.from_single_file
] for loading.ckpt
files - [
~loaders.IPAdapterMixin.load_ip_adapter
] for loading IP Adapters
PARAMETER | DESCRIPTION |
---|---|
vae |
Variational Auto-Encoder (VAE) model to encode and decode images to and from latent representations.
TYPE:
|
text_encoder |
Frozen text-encoder (clip-vit-large-patch14).
TYPE:
|
tokenizer |
A
TYPE:
|
unet |
A
TYPE:
|
scheduler |
A scheduler to be used in combination with
TYPE:
|
safety_checker |
Classification module that estimates whether generated images could be considered offensive or harmful. Please refer to the model card for more details about a model's potential harms.
TYPE:
|
feature_extractor |
A
TYPE:
|
Source code in mindone/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion.py
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|
mindone.diffusers.StableDiffusionPipeline.__call__(prompt=None, height=None, width=None, num_inference_steps=50, timesteps=None, sigmas=None, guidance_scale=7.5, negative_prompt=None, num_images_per_prompt=1, eta=0.0, generator=None, latents=None, prompt_embeds=None, negative_prompt_embeds=None, ip_adapter_image=None, ip_adapter_image_embeds=None, output_type='pil', return_dict=False, cross_attention_kwargs=None, guidance_rescale=0.0, clip_skip=None, callback_on_step_end=None, callback_on_step_end_tensor_inputs=['latents'], **kwargs)
¶
The call function to the pipeline for generation.
PARAMETER | DESCRIPTION |
---|---|
prompt |
The prompt or prompts to guide image generation. If not defined, you need to pass
TYPE:
|
height |
The height in pixels of the generated image.
TYPE:
|
width |
The width in pixels of the generated image.
TYPE:
|
num_inference_steps |
The number of denoising steps. More denoising steps usually lead to a higher quality image at the expense of slower inference.
TYPE:
|
timesteps |
Custom timesteps to use for the denoising process with schedulers which support a
TYPE:
|
sigmas |
Custom sigmas to use for the denoising process with schedulers which support a
TYPE:
|
guidance_scale |
A higher guidance scale value encourages the model to generate images closely linked to the text
TYPE:
|
negative_prompt |
The prompt or prompts to guide what to not include in image generation. If not defined, you need to
pass
TYPE:
|
num_images_per_prompt |
The number of images to generate per prompt.
TYPE:
|
eta |
Corresponds to parameter eta (η) from the DDIM paper. Only applies
to the [
TYPE:
|
generator |
A
TYPE:
|
latents |
Pre-generated noisy latents sampled from a Gaussian distribution, to be used as inputs for image
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents
tensor is generated by sampling using the supplied random
TYPE:
|
prompt_embeds |
Pre-generated text embeddings. Can be used to easily tweak text inputs (prompt weighting). If not
provided, text embeddings are generated from the
TYPE:
|
negative_prompt_embeds |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs (prompt weighting). If
not provided,
TYPE:
|
ip_adapter_image |
(
TYPE:
|
ip_adapter_image_embeds |
Pre-generated image embeddings for IP-Adapter. It should be a list of length same as number of
IP-adapters. Each element should be a tensor of shape
TYPE:
|
output_type |
The output format of the generated image. Choose between
TYPE:
|
return_dict |
Whether or not to return a [
TYPE:
|
cross_attention_kwargs |
A kwargs dictionary that if specified is passed along to the [
TYPE:
|
guidance_rescale |
Guidance rescale factor from Common Diffusion Noise Schedules and Sample Steps are Flawed. Guidance rescale factor should fix overexposure when using zero terminal SNR.
TYPE:
|
clip_skip |
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings.
TYPE:
|
callback_on_step_end |
A function or a subclass of
TYPE:
|
callback_on_step_end_tensor_inputs |
The list of tensor inputs for the
TYPE:
|
RETURNS | DESCRIPTION |
---|---|
[ |
Source code in mindone/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion.py
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|
mindone.diffusers.StableDiffusionPipeline.encode_prompt(prompt, num_images_per_prompt, do_classifier_free_guidance, negative_prompt=None, prompt_embeds=None, negative_prompt_embeds=None, lora_scale=None, clip_skip=None)
¶
Encodes the prompt into text encoder hidden states.
PARAMETER | DESCRIPTION |
---|---|
prompt |
prompt to be encoded
TYPE:
|
num_images_per_prompt |
number of images that should be generated per prompt
TYPE:
|
do_classifier_free_guidance |
whether to use classifier free guidance or not
TYPE:
|
negative_prompt |
The prompt or prompts not to guide the image generation. If not defined, one has to pass
TYPE:
|
prompt_embeds |
Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not
provided, text embeddings will be generated from
TYPE:
|
negative_prompt_embeds |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, negative_prompt_embeds will be generated from
TYPE:
|
lora_scale |
A LoRA scale that will be applied to all LoRA layers of the text encoder if LoRA layers are loaded.
TYPE:
|
clip_skip |
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings.
TYPE:
|
Source code in mindone/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion.py
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|
mindone.diffusers.StableDiffusionPipeline.get_guidance_scale_embedding(w, embedding_dim=512, dtype=ms.float32)
¶
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
PARAMETER | DESCRIPTION |
---|---|
w |
Generate embedding vectors with a specified guidance scale to subsequently enrich timestep embeddings.
TYPE:
|
embedding_dim |
Dimension of the embeddings to generate.
TYPE:
|
dtype |
Data type of the generated embeddings.
TYPE:
|
RETURNS | DESCRIPTION |
---|---|
Tensor
|
|
Source code in mindone/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion.py
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|
mindone.diffusers.pipelines.stable_diffusion.StableDiffusionPipelineOutput
dataclass
¶
Bases: BaseOutput
Output class for Stable Diffusion pipelines.
Source code in mindone/diffusers/pipelines/stable_diffusion/pipeline_output.py
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