Stable Diffusion XL¶
Stable Diffusion XL (SDXL) was proposed in SDXL: Improving Latent Diffusion Models for High-Resolution Image Synthesis by Dustin Podell, Zion English, Kyle Lacey, Andreas Blattmann, Tim Dockhorn, Jonas Müller, Joe Penna, and Robin Rombach.
The abstract from the paper is:
We present SDXL, a latent diffusion model for text-to-image synthesis. Compared to previous versions of Stable Diffusion, SDXL leverages a three times larger UNet backbone: The increase of model parameters is mainly due to more attention blocks and a larger cross-attention context as SDXL uses a second text encoder. We design multiple novel conditioning schemes and train SDXL on multiple aspect ratios. We also introduce a refinement model which is used to improve the visual fidelity of samples generated by SDXL using a post-hoc image-to-image technique. We demonstrate that SDXL shows drastically improved performance compared the previous versions of Stable Diffusion and achieves results competitive with those of black-box state-of-the-art image generators.
Tips¶
- Using SDXL with a DPM++ scheduler for less than 50 steps is known to produce visual artifacts because the solver becomes numerically unstable. To fix this issue, take a look at this PR which recommends for ODE/SDE solvers:
- set
use_karras_sigmas=True
orlu_lambdas=True
to improve image quality - set
euler_at_final=True
if you're using a solver with uniform step sizes (DPM++2M or DPM++2M SDE)
- set
- Most SDXL checkpoints work best with an image size of 1024x1024. Image sizes of 768x768 and 512x512 are also supported, but the results aren't as good. Anything below 512x512 is not recommended and likely won't be for default checkpoints like stabilityai/stable-diffusion-xl-base-1.0.
- SDXL can pass a different prompt for each of the text encoders it was trained on. We can even pass different parts of the same prompt to the text encoders.
- SDXL output images can be improved by making use of a refiner model in an image-to-image setting.
- SDXL offers
negative_original_size
,negative_crops_coords_top_left
, andnegative_target_size
to negatively condition the model on image resolution and cropping parameters.
Tip
To learn how to use SDXL for various tasks, how to optimize performance, and other usage examples, take a look at the Stable Diffusion XL guide.
Check out the Stability AI Hub organization for the official base and refiner model checkpoints!
mindone.diffusers.StableDiffusionXLPipeline
¶
Bases: DiffusionPipeline
, StableDiffusionMixin
, FromSingleFileMixin
, IPAdapterMixin
, StableDiffusionXLLoraLoaderMixin
, TextualInversionLoaderMixin
Pipeline for text-to-image generation using Stable Diffusion XL.
This model inherits from [DiffusionPipeline
]. Check the superclass documentation for the generic methods the
library implements for all the pipelines (such as downloading or saving, running on a particular device, etc.)
The pipeline also inherits the following loading methods
- [
~loaders.TextualInversionLoaderMixin.load_textual_inversion
] for loading textual inversion embeddings - [
~loaders.FromSingleFileMixin.from_single_file
] for loading.ckpt
files - [
~loaders.StableDiffusionXLLoraLoaderMixin.load_lora_weights
] for loading LoRA weights - [
~loaders.StableDiffusionXLLoraLoaderMixin.save_lora_weights
] for saving LoRA weights - [
~loaders.IPAdapterMixin.load_ip_adapter
] for loading IP Adapters
PARAMETER | DESCRIPTION |
---|---|
vae |
Variational Auto-Encoder (VAE) Model to encode and decode images to and from latent representations.
TYPE:
|
text_encoder |
Frozen text-encoder. Stable Diffusion XL uses the text portion of CLIP, specifically the clip-vit-large-patch14 variant.
TYPE:
|
text_encoder_2 |
Second frozen text-encoder. Stable Diffusion XL uses the text and pool portion of CLIP, specifically the laion/CLIP-ViT-bigG-14-laion2B-39B-b160k variant.
TYPE:
|
tokenizer |
Tokenizer of class CLIPTokenizer.
TYPE:
|
tokenizer_2 |
Second Tokenizer of class CLIPTokenizer.
TYPE:
|
unet |
Conditional U-Net architecture to denoise the encoded image latents.
TYPE:
|
scheduler |
A scheduler to be used in combination with
TYPE:
|
force_zeros_for_empty_prompt |
Whether the negative prompt embeddings shall be forced to always be set to 0. Also see the config of
TYPE:
|
add_watermarker |
Whether to use the invisible_watermark library to watermark output images. If not defined, it will default to True if the package is installed, otherwise no watermarker will be used.
TYPE:
|
Source code in mindone/diffusers/pipelines/stable_diffusion_xl/pipeline_stable_diffusion_xl.py
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|
mindone.diffusers.StableDiffusionXLPipeline.__call__(prompt=None, prompt_2=None, height=None, width=None, num_inference_steps=50, timesteps=None, sigmas=None, denoising_end=None, guidance_scale=5.0, negative_prompt=None, negative_prompt_2=None, num_images_per_prompt=1, eta=0.0, generator=None, latents=None, prompt_embeds=None, negative_prompt_embeds=None, pooled_prompt_embeds=None, negative_pooled_prompt_embeds=None, ip_adapter_image=None, ip_adapter_image_embeds=None, output_type='pil', return_dict=False, cross_attention_kwargs=None, guidance_rescale=0.0, original_size=None, crops_coords_top_left=(0, 0), target_size=None, negative_original_size=None, negative_crops_coords_top_left=(0, 0), negative_target_size=None, clip_skip=None, callback_on_step_end=None, callback_on_step_end_tensor_inputs=['latents'], **kwargs)
¶
Function invoked when calling the pipeline for generation.
PARAMETER | DESCRIPTION |
---|---|
prompt |
The prompt or prompts to guide the image generation. If not defined, one has to pass
TYPE:
|
prompt_2 |
The prompt or prompts to be sent to the
TYPE:
|
height |
The height in pixels of the generated image. This is set to 1024 by default for the best results. Anything below 512 pixels won't work well for stabilityai/stable-diffusion-xl-base-1.0 and checkpoints that are not specifically fine-tuned on low resolutions.
TYPE:
|
width |
The width in pixels of the generated image. This is set to 1024 by default for the best results. Anything below 512 pixels won't work well for stabilityai/stable-diffusion-xl-base-1.0 and checkpoints that are not specifically fine-tuned on low resolutions.
TYPE:
|
num_inference_steps |
The number of denoising steps. More denoising steps usually lead to a higher quality image at the expense of slower inference.
TYPE:
|
timesteps |
Custom timesteps to use for the denoising process with schedulers which support a
TYPE:
|
sigmas |
Custom sigmas to use for the denoising process with schedulers which support a
TYPE:
|
denoising_end |
When specified, determines the fraction (between 0.0 and 1.0) of the total denoising process to be completed before it is intentionally prematurely terminated. As a result, the returned sample will still retain a substantial amount of noise as determined by the discrete timesteps selected by the scheduler. The denoising_end parameter should ideally be utilized when this pipeline forms a part of a "Mixture of Denoisers" multi-pipeline setup, as elaborated in Refining the Image Output
TYPE:
|
guidance_scale |
Guidance scale as defined in Classifier-Free Diffusion Guidance.
TYPE:
|
negative_prompt |
The prompt or prompts not to guide the image generation. If not defined, one has to pass
TYPE:
|
negative_prompt_2 |
The prompt or prompts not to guide the image generation to be sent to
TYPE:
|
num_images_per_prompt |
The number of images to generate per prompt.
TYPE:
|
eta |
Corresponds to parameter eta (η) in the DDIM paper: https://arxiv.org/abs/2010.02502. Only applies to
[
TYPE:
|
generator |
One or a list of numpy generator(s) to make generation deterministic.
TYPE:
|
latents |
Pre-generated noisy latents, sampled from a Gaussian distribution, to be used as inputs for image
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents
tensor will ge generated by sampling using the supplied random
TYPE:
|
prompt_embeds |
Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not
provided, text embeddings will be generated from
TYPE:
|
negative_prompt_embeds |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, negative_prompt_embeds will be generated from
TYPE:
|
pooled_prompt_embeds |
Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting.
If not provided, pooled text embeddings will be generated from
TYPE:
|
negative_pooled_prompt_embeds |
Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, pooled negative_prompt_embeds will be generated from
TYPE:
|
ip_adapter_image |
(
TYPE:
|
ip_adapter_image_embeds |
Pre-generated image embeddings for IP-Adapter. It should be a list of length same as number of IP-adapters.
Each element should be a tensor of shape
TYPE:
|
output_type |
The output format of the generate image. Choose between
PIL:
TYPE:
|
return_dict |
Whether or not to return a [
TYPE:
|
cross_attention_kwargs |
A kwargs dictionary that if specified is passed along to the
TYPE:
|
guidance_rescale |
Guidance rescale factor proposed by Common Diffusion Noise Schedules and Sample Steps are
Flawed
TYPE:
|
original_size |
If
TYPE:
|
crops_coords_top_left |
TYPE:
|
target_size |
For most cases,
TYPE:
|
negative_original_size |
To negatively condition the generation process based on a specific image resolution. Part of SDXL's micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. For more information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208.
TYPE:
|
negative_crops_coords_top_left |
To negatively condition the generation process based on a specific crop coordinates. Part of SDXL's micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. For more information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208.
TYPE:
|
negative_target_size |
To negatively condition the generation process based on a target image resolution. It should be as same
as the
TYPE:
|
callback_on_step_end |
A function or a subclass of
TYPE:
|
callback_on_step_end_tensor_inputs |
The list of tensor inputs for the
TYPE:
|
RETURNS | DESCRIPTION |
---|---|
[ |
|
[ |
|
|
Source code in mindone/diffusers/pipelines/stable_diffusion_xl/pipeline_stable_diffusion_xl.py
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|
mindone.diffusers.StableDiffusionXLPipeline.encode_prompt(prompt, prompt_2=None, num_images_per_prompt=1, do_classifier_free_guidance=True, negative_prompt=None, negative_prompt_2=None, prompt_embeds=None, negative_prompt_embeds=None, pooled_prompt_embeds=None, negative_pooled_prompt_embeds=None, lora_scale=None, clip_skip=None)
¶
Encodes the prompt into text encoder hidden states.
PARAMETER | DESCRIPTION |
---|---|
prompt |
prompt to be encoded
TYPE:
|
prompt_2 |
The prompt or prompts to be sent to the
TYPE:
|
num_images_per_prompt |
number of images that should be generated per prompt
TYPE:
|
do_classifier_free_guidance |
whether to use classifier free guidance or not
TYPE:
|
negative_prompt |
The prompt or prompts not to guide the image generation. If not defined, one has to pass
TYPE:
|
negative_prompt_2 |
The prompt or prompts not to guide the image generation to be sent to
TYPE:
|
prompt_embeds |
Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not
provided, text embeddings will be generated from
TYPE:
|
negative_prompt_embeds |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, negative_prompt_embeds will be generated from
TYPE:
|
pooled_prompt_embeds |
Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting.
If not provided, pooled text embeddings will be generated from
TYPE:
|
negative_pooled_prompt_embeds |
Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, pooled negative_prompt_embeds will be generated from
TYPE:
|
lora_scale |
A lora scale that will be applied to all LoRA layers of the text encoder if LoRA layers are loaded.
TYPE:
|
clip_skip |
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings.
TYPE:
|
Source code in mindone/diffusers/pipelines/stable_diffusion_xl/pipeline_stable_diffusion_xl.py
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|
mindone.diffusers.StableDiffusionXLPipeline.get_guidance_scale_embedding(w, embedding_dim=512, dtype=ms.float32)
¶
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
PARAMETER | DESCRIPTION |
---|---|
w |
Generate embedding vectors with a specified guidance scale to subsequently enrich timestep embeddings.
TYPE:
|
embedding_dim |
Dimension of the embeddings to generate.
TYPE:
|
dtype |
Data type of the generated embeddings.
TYPE:
|
RETURNS | DESCRIPTION |
---|---|
Tensor
|
|
Source code in mindone/diffusers/pipelines/stable_diffusion_xl/pipeline_stable_diffusion_xl.py
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|
mindone.diffusers.StableDiffusionXLImg2ImgPipeline
¶
Bases: DiffusionPipeline
, StableDiffusionMixin
, FromSingleFileMixin
, IPAdapterMixin
, StableDiffusionXLLoraLoaderMixin
, TextualInversionLoaderMixin
Pipeline for text-to-image generation using Stable Diffusion XL.
This model inherits from [DiffusionPipeline
]. Check the superclass documentation for the generic methods the
library implements for all the pipelines (such as downloading or saving, running on a particular device, etc.)
The pipeline also inherits the following loading methods
- [
~loaders.TextualInversionLoaderMixin.load_textual_inversion
] for loading textual inversion embeddings - [
~loaders.FromSingleFileMixin.from_single_file
] for loading.ckpt
files - [
~loaders.StableDiffusionXLLoraLoaderMixin.load_lora_weights
] for loading LoRA weights - [
~loaders.StableDiffusionXLLoraLoaderMixin.save_lora_weights
] for saving LoRA weights - [
~loaders.IPAdapterMixin.load_ip_adapter
] for loading IP Adapters
PARAMETER | DESCRIPTION |
---|---|
vae |
Variational Auto-Encoder (VAE) Model to encode and decode images to and from latent representations.
TYPE:
|
text_encoder |
Frozen text-encoder. Stable Diffusion XL uses the text portion of CLIP, specifically the clip-vit-large-patch14 variant.
TYPE:
|
text_encoder_2 |
Second frozen text-encoder. Stable Diffusion XL uses the text and pool portion of CLIP, specifically the laion/CLIP-ViT-bigG-14-laion2B-39B-b160k variant.
TYPE:
|
tokenizer |
Tokenizer of class CLIPTokenizer.
TYPE:
|
tokenizer_2 |
Second Tokenizer of class CLIPTokenizer.
TYPE:
|
unet |
Conditional U-Net architecture to denoise the encoded image latents.
TYPE:
|
scheduler |
A scheduler to be used in combination with
TYPE:
|
requires_aesthetics_score |
Whether the
TYPE:
|
force_zeros_for_empty_prompt |
Whether the negative prompt embeddings shall be forced to always be set to 0. Also see the config of
TYPE:
|
add_watermarker |
Whether to use the invisible_watermark library to watermark output images. If not defined, it will default to True if the package is installed, otherwise no watermarker will be used.
TYPE:
|
Source code in mindone/diffusers/pipelines/stable_diffusion_xl/pipeline_stable_diffusion_xl_img2img.py
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mindone.diffusers.StableDiffusionXLImg2ImgPipeline.__call__(prompt=None, prompt_2=None, image=None, strength=0.3, num_inference_steps=50, timesteps=None, sigmas=None, denoising_start=None, denoising_end=None, guidance_scale=5.0, negative_prompt=None, negative_prompt_2=None, num_images_per_prompt=1, eta=0.0, generator=None, latents=None, prompt_embeds=None, negative_prompt_embeds=None, pooled_prompt_embeds=None, negative_pooled_prompt_embeds=None, ip_adapter_image=None, ip_adapter_image_embeds=None, output_type='pil', return_dict=False, cross_attention_kwargs=None, guidance_rescale=0.0, original_size=None, crops_coords_top_left=(0, 0), target_size=None, negative_original_size=None, negative_crops_coords_top_left=(0, 0), negative_target_size=None, aesthetic_score=6.0, negative_aesthetic_score=2.5, clip_skip=None, callback_on_step_end=None, callback_on_step_end_tensor_inputs=['latents'], **kwargs)
¶
Function invoked when calling the pipeline for generation.
PARAMETER | DESCRIPTION |
---|---|
prompt |
The prompt or prompts to guide the image generation. If not defined, one has to pass
TYPE:
|
prompt_2 |
The prompt or prompts to be sent to the
TYPE:
|
image |
The image(s) to modify with the pipeline.
TYPE:
|
strength |
Conceptually, indicates how much to transform the reference
TYPE:
|
num_inference_steps |
The number of denoising steps. More denoising steps usually lead to a higher quality image at the expense of slower inference.
TYPE:
|
timesteps |
Custom timesteps to use for the denoising process with schedulers which support a
TYPE:
|
sigmas |
Custom sigmas to use for the denoising process with schedulers which support a
TYPE:
|
denoising_start |
When specified, indicates the fraction (between 0.0 and 1.0) of the total denoising process to be
bypassed before it is initiated. Consequently, the initial part of the denoising process is skipped and
it is assumed that the passed
TYPE:
|
denoising_end |
When specified, determines the fraction (between 0.0 and 1.0) of the total denoising process to be
completed before it is intentionally prematurely terminated. As a result, the returned sample will
still retain a substantial amount of noise (ca. final 20% of timesteps still needed) and should be
denoised by a successor pipeline that has
TYPE:
|
guidance_scale |
Guidance scale as defined in Classifier-Free Diffusion Guidance.
TYPE:
|
negative_prompt |
The prompt or prompts not to guide the image generation. If not defined, one has to pass
TYPE:
|
negative_prompt_2 |
The prompt or prompts not to guide the image generation to be sent to
TYPE:
|
num_images_per_prompt |
The number of images to generate per prompt.
TYPE:
|
eta |
Corresponds to parameter eta (η) in the DDIM paper: https://arxiv.org/abs/2010.02502. Only applies to
[
TYPE:
|
generator |
One or a list of numpy generator(s) to make generation deterministic.
TYPE:
|
latents |
Pre-generated noisy latents, sampled from a Gaussian distribution, to be used as inputs for image
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents
tensor will ge generated by sampling using the supplied random
TYPE:
|
prompt_embeds |
Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not
provided, text embeddings will be generated from
TYPE:
|
negative_prompt_embeds |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, negative_prompt_embeds will be generated from
TYPE:
|
pooled_prompt_embeds |
Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting.
If not provided, pooled text embeddings will be generated from
TYPE:
|
negative_pooled_prompt_embeds |
Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, pooled negative_prompt_embeds will be generated from
TYPE:
|
ip_adapter_image |
(
TYPE:
|
ip_adapter_image_embeds |
Pre-generated image embeddings for IP-Adapter. It should be a list of length same as number of IP-adapters.
Each element should be a tensor of shape
TYPE:
|
output_type |
The output format of the generate image. Choose between
PIL:
TYPE:
|
return_dict |
Whether or not to return a [
TYPE:
|
cross_attention_kwargs |
A kwargs dictionary that if specified is passed along to the
TYPE:
|
guidance_rescale |
Guidance rescale factor proposed by Common Diffusion Noise Schedules and Sample Steps are
Flawed
TYPE:
|
original_size |
If
TYPE:
|
crops_coords_top_left |
TYPE:
|
target_size |
For most cases,
TYPE:
|
negative_original_size |
To negatively condition the generation process based on a specific image resolution. Part of SDXL's micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. For more information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208.
TYPE:
|
negative_crops_coords_top_left |
To negatively condition the generation process based on a specific crop coordinates. Part of SDXL's micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. For more information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208.
TYPE:
|
negative_target_size |
To negatively condition the generation process based on a target image resolution. It should be as same
as the
TYPE:
|
aesthetic_score |
Used to simulate an aesthetic score of the generated image by influencing the positive text condition. Part of SDXL's micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952.
TYPE:
|
negative_aesthetic_score |
Part of SDXL's micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. Can be used to simulate an aesthetic score of the generated image by influencing the negative text condition.
TYPE:
|
clip_skip |
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings.
TYPE:
|
callback_on_step_end |
A function or a subclass of
TYPE:
|
callback_on_step_end_tensor_inputs |
The list of tensor inputs for the
TYPE:
|
RETURNS | DESCRIPTION |
---|---|
[ |
|
[ |
|
`tuple. When returning a tuple, the first element is a list with the generated images. |
Source code in mindone/diffusers/pipelines/stable_diffusion_xl/pipeline_stable_diffusion_xl_img2img.py
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mindone.diffusers.StableDiffusionXLImg2ImgPipeline.encode_prompt(prompt, prompt_2=None, num_images_per_prompt=1, do_classifier_free_guidance=True, negative_prompt=None, negative_prompt_2=None, prompt_embeds=None, negative_prompt_embeds=None, pooled_prompt_embeds=None, negative_pooled_prompt_embeds=None, lora_scale=None, clip_skip=None)
¶
Encodes the prompt into text encoder hidden states.
PARAMETER | DESCRIPTION |
---|---|
prompt |
prompt to be encoded
TYPE:
|
prompt_2 |
The prompt or prompts to be sent to the
TYPE:
|
num_images_per_prompt |
number of images that should be generated per prompt
TYPE:
|
do_classifier_free_guidance |
whether to use classifier free guidance or not
TYPE:
|
negative_prompt |
The prompt or prompts not to guide the image generation. If not defined, one has to pass
TYPE:
|
negative_prompt_2 |
The prompt or prompts not to guide the image generation to be sent to
TYPE:
|
prompt_embeds |
Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not
provided, text embeddings will be generated from
TYPE:
|
negative_prompt_embeds |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, negative_prompt_embeds will be generated from
TYPE:
|
pooled_prompt_embeds |
Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting.
If not provided, pooled text embeddings will be generated from
TYPE:
|
negative_pooled_prompt_embeds |
Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, pooled negative_prompt_embeds will be generated from
TYPE:
|
lora_scale |
A lora scale that will be applied to all LoRA layers of the text encoder if LoRA layers are loaded.
TYPE:
|
clip_skip |
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings.
TYPE:
|
Source code in mindone/diffusers/pipelines/stable_diffusion_xl/pipeline_stable_diffusion_xl_img2img.py
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|
mindone.diffusers.StableDiffusionXLImg2ImgPipeline.get_guidance_scale_embedding(w, embedding_dim=512, dtype=ms.float32)
¶
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
PARAMETER | DESCRIPTION |
---|---|
w |
Generate embedding vectors with a specified guidance scale to subsequently enrich timestep embeddings.
TYPE:
|
embedding_dim |
Dimension of the embeddings to generate.
TYPE:
|
dtype |
Data type of the generated embeddings.
TYPE:
|
RETURNS | DESCRIPTION |
---|---|
Tensor
|
|
Source code in mindone/diffusers/pipelines/stable_diffusion_xl/pipeline_stable_diffusion_xl_img2img.py
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|
mindone.diffusers.StableDiffusionXLInpaintPipeline
¶
Bases: DiffusionPipeline
, StableDiffusionMixin
, FromSingleFileMixin
, IPAdapterMixin
, StableDiffusionXLLoraLoaderMixin
, TextualInversionLoaderMixin
Pipeline for text-to-image generation using Stable Diffusion XL.
This model inherits from [DiffusionPipeline
]. Check the superclass documentation for the generic methods the
library implements for all the pipelines (such as downloading or saving, running on a particular device, etc.)
The pipeline also inherits the following loading methods
- [
~loaders.TextualInversionLoaderMixin.load_textual_inversion
] for loading textual inversion embeddings - [
~loaders.FromSingleFileMixin.from_single_file
] for loading.ckpt
files - [
~loaders.StableDiffusionXLLoraLoaderMixin.load_lora_weights
] for loading LoRA weights - [
~loaders.StableDiffusionXLLoraLoaderMixin.save_lora_weights
] for saving LoRA weights - [
~loaders.IPAdapterMixin.load_ip_adapter
] for loading IP Adapters
PARAMETER | DESCRIPTION |
---|---|
vae |
Variational Auto-Encoder (VAE) Model to encode and decode images to and from latent representations.
TYPE:
|
text_encoder |
Frozen text-encoder. Stable Diffusion XL uses the text portion of CLIP, specifically the clip-vit-large-patch14 variant.
TYPE:
|
text_encoder_2 |
Second frozen text-encoder. Stable Diffusion XL uses the text and pool portion of CLIP, specifically the laion/CLIP-ViT-bigG-14-laion2B-39B-b160k variant.
TYPE:
|
tokenizer |
Tokenizer of class CLIPTokenizer.
TYPE:
|
tokenizer_2 |
Second Tokenizer of class CLIPTokenizer.
TYPE:
|
unet |
Conditional U-Net architecture to denoise the encoded image latents.
TYPE:
|
scheduler |
A scheduler to be used in combination with
TYPE:
|
requires_aesthetics_score |
Whether the
TYPE:
|
force_zeros_for_empty_prompt |
Whether the negative prompt embeddings shall be forced to always be set to 0. Also see the config of
TYPE:
|
add_watermarker |
Whether to use the invisible_watermark library to watermark output images. If not defined, it will default to True if the package is installed, otherwise no watermarker will be used.
TYPE:
|
Source code in mindone/diffusers/pipelines/stable_diffusion_xl/pipeline_stable_diffusion_xl_inpaint.py
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mindone.diffusers.StableDiffusionXLInpaintPipeline.__call__(prompt=None, prompt_2=None, image=None, mask_image=None, masked_image_latents=None, height=None, width=None, padding_mask_crop=None, strength=0.9999, num_inference_steps=50, timesteps=None, sigmas=None, denoising_start=None, denoising_end=None, guidance_scale=7.5, negative_prompt=None, negative_prompt_2=None, num_images_per_prompt=1, eta=0.0, generator=None, latents=None, prompt_embeds=None, negative_prompt_embeds=None, pooled_prompt_embeds=None, negative_pooled_prompt_embeds=None, ip_adapter_image=None, ip_adapter_image_embeds=None, output_type='pil', return_dict=False, cross_attention_kwargs=None, guidance_rescale=0.0, original_size=None, crops_coords_top_left=(0, 0), target_size=None, negative_original_size=None, negative_crops_coords_top_left=(0, 0), negative_target_size=None, aesthetic_score=6.0, negative_aesthetic_score=2.5, clip_skip=None, callback_on_step_end=None, callback_on_step_end_tensor_inputs=['latents'], **kwargs)
¶
Function invoked when calling the pipeline for generation.
PARAMETER | DESCRIPTION |
---|---|
prompt |
The prompt or prompts to guide the image generation. If not defined, one has to pass
TYPE:
|
prompt_2 |
The prompt or prompts to be sent to the
TYPE:
|
image |
TYPE:
|
mask_image |
TYPE:
|
height |
The height in pixels of the generated image. This is set to 1024 by default for the best results. Anything below 512 pixels won't work well for stabilityai/stable-diffusion-xl-base-1.0 and checkpoints that are not specifically fine-tuned on low resolutions.
TYPE:
|
width |
The width in pixels of the generated image. This is set to 1024 by default for the best results. Anything below 512 pixels won't work well for stabilityai/stable-diffusion-xl-base-1.0 and checkpoints that are not specifically fine-tuned on low resolutions.
TYPE:
|
padding_mask_crop |
The size of margin in the crop to be applied to the image and masking. If
TYPE:
|
strength |
Conceptually, indicates how much to transform the masked portion of the reference
TYPE:
|
num_inference_steps |
The number of denoising steps. More denoising steps usually lead to a higher quality image at the expense of slower inference.
TYPE:
|
timesteps |
Custom timesteps to use for the denoising process with schedulers which support a
TYPE:
|
sigmas |
Custom sigmas to use for the denoising process with schedulers which support a
TYPE:
|
denoising_start |
When specified, indicates the fraction (between 0.0 and 1.0) of the total denoising process to be
bypassed before it is initiated. Consequently, the initial part of the denoising process is skipped and
it is assumed that the passed
TYPE:
|
denoising_end |
When specified, determines the fraction (between 0.0 and 1.0) of the total denoising process to be
completed before it is intentionally prematurely terminated. As a result, the returned sample will
still retain a substantial amount of noise (ca. final 20% of timesteps still needed) and should be
denoised by a successor pipeline that has
TYPE:
|
guidance_scale |
Guidance scale as defined in Classifier-Free Diffusion Guidance.
TYPE:
|
negative_prompt |
The prompt or prompts not to guide the image generation. If not defined, one has to pass
TYPE:
|
negative_prompt_2 |
The prompt or prompts not to guide the image generation to be sent to
TYPE:
|
prompt_embeds |
Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not
provided, text embeddings will be generated from
TYPE:
|
negative_prompt_embeds |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, negative_prompt_embeds will be generated from
TYPE:
|
pooled_prompt_embeds |
Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting.
If not provided, pooled text embeddings will be generated from
TYPE:
|
negative_pooled_prompt_embeds |
Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, pooled negative_prompt_embeds will be generated from
TYPE:
|
ip_adapter_image |
(
TYPE:
|
ip_adapter_image_embeds |
Pre-generated image embeddings for IP-Adapter. It should be a list of length same as number of IP-adapters.
Each element should be a tensor of shape
TYPE:
|
num_images_per_prompt |
The number of images to generate per prompt.
TYPE:
|
eta |
Corresponds to parameter eta (??) in the DDIM paper: https://arxiv.org/abs/2010.02502. Only applies to
[
TYPE:
|
generator |
One or a list of numpy generator(s) to make generation deterministic.
TYPE:
|
latents |
Pre-generated noisy latents, sampled from a Gaussian distribution, to be used as inputs for image
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents
tensor will ge generated by sampling using the supplied random
TYPE:
|
output_type |
The output format of the generate image. Choose between
PIL:
TYPE:
|
return_dict |
Whether or not to return a [
TYPE:
|
cross_attention_kwargs |
A kwargs dictionary that if specified is passed along to the
TYPE:
|
original_size |
If
TYPE:
|
crops_coords_top_left |
TYPE:
|
target_size |
For most cases,
TYPE:
|
negative_original_size |
To negatively condition the generation process based on a specific image resolution. Part of SDXL's micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. For more information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208.
TYPE:
|
negative_crops_coords_top_left |
To negatively condition the generation process based on a specific crop coordinates. Part of SDXL's micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. For more information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208.
TYPE:
|
negative_target_size |
To negatively condition the generation process based on a target image resolution. It should be as same
as the
TYPE:
|
aesthetic_score |
Used to simulate an aesthetic score of the generated image by influencing the positive text condition. Part of SDXL's micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952.
TYPE:
|
negative_aesthetic_score |
Part of SDXL's micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. Can be used to simulate an aesthetic score of the generated image by influencing the negative text condition.
TYPE:
|
clip_skip |
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings.
TYPE:
|
callback_on_step_end |
A function or a subclass of
TYPE:
|
callback_on_step_end_tensor_inputs |
The list of tensor inputs for the
TYPE:
|
RETURNS | DESCRIPTION |
---|---|
[ |
|
[ |
|
|
Source code in mindone/diffusers/pipelines/stable_diffusion_xl/pipeline_stable_diffusion_xl_inpaint.py
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mindone.diffusers.StableDiffusionXLInpaintPipeline.encode_prompt(prompt, prompt_2=None, num_images_per_prompt=1, do_classifier_free_guidance=True, negative_prompt=None, negative_prompt_2=None, prompt_embeds=None, negative_prompt_embeds=None, pooled_prompt_embeds=None, negative_pooled_prompt_embeds=None, lora_scale=None, clip_skip=None)
¶
Encodes the prompt into text encoder hidden states.
PARAMETER | DESCRIPTION |
---|---|
prompt |
prompt to be encoded
TYPE:
|
prompt_2 |
The prompt or prompts to be sent to the
TYPE:
|
num_images_per_prompt |
number of images that should be generated per prompt
TYPE:
|
do_classifier_free_guidance |
whether to use classifier free guidance or not
TYPE:
|
negative_prompt |
The prompt or prompts not to guide the image generation. If not defined, one has to pass
TYPE:
|
negative_prompt_2 |
The prompt or prompts not to guide the image generation to be sent to
TYPE:
|
prompt_embeds |
Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not
provided, text embeddings will be generated from
TYPE:
|
negative_prompt_embeds |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, negative_prompt_embeds will be generated from
TYPE:
|
pooled_prompt_embeds |
Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting.
If not provided, pooled text embeddings will be generated from
TYPE:
|
negative_pooled_prompt_embeds |
Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, pooled negative_prompt_embeds will be generated from
TYPE:
|
lora_scale |
A lora scale that will be applied to all LoRA layers of the text encoder if LoRA layers are loaded.
TYPE:
|
clip_skip |
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings.
TYPE:
|
Source code in mindone/diffusers/pipelines/stable_diffusion_xl/pipeline_stable_diffusion_xl_inpaint.py
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mindone.diffusers.StableDiffusionXLInpaintPipeline.get_guidance_scale_embedding(w, embedding_dim=512, dtype=ms.float32)
¶
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
PARAMETER | DESCRIPTION |
---|---|
w |
Generate embedding vectors with a specified guidance scale to subsequently enrich timestep embeddings.
TYPE:
|
embedding_dim |
Dimension of the embeddings to generate.
TYPE:
|
dtype |
Data type of the generated embeddings.
TYPE:
|
RETURNS | DESCRIPTION |
---|---|
Tensor
|
|
Source code in mindone/diffusers/pipelines/stable_diffusion_xl/pipeline_stable_diffusion_xl_inpaint.py
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