GLIGEN (Grounded Language-to-Image Generation)¶
The GLIGEN model was created by researchers and engineers from University of Wisconsin-Madison, Columbia University, and Microsoft. The StableDiffusionGLIGENPipeline
and StableDiffusionGLIGENTextImagePipeline
can generate photorealistic images conditioned on grounding inputs. Along with text and bounding boxes with StableDiffusionGLIGENPipeline
, if input images are given, StableDiffusionGLIGENTextImagePipeline
can insert objects described by text at the region defined by bounding boxes. Otherwise, it'll generate an image described by the caption/prompt and insert objects described by text at the region defined by bounding boxes. It's trained on COCO2014D and COCO2014CD datasets, and the model uses a frozen CLIP ViT-L/14 text encoder to condition itself on grounding inputs.
The abstract from the paper is:
Large-scale text-to-image diffusion models have made amazing advances. However, the status quo is to use text input alone, which can impede controllability. In this work, we propose GLIGEN, Grounded-Language-to-Image Generation, a novel approach that builds upon and extends the functionality of existing pre-trained text-to-image diffusion models by enabling them to also be conditioned on grounding inputs. To preserve the vast concept knowledge of the pre-trained model, we freeze all of its weights and inject the grounding information into new trainable layers via a gated mechanism. Our model achieves open-world grounded text2img generation with caption and bounding box condition inputs, and the grounding ability generalizes well to novel spatial configurations and concepts. GLIGEN’s zeroshot performance on COCO and LVIS outperforms existing supervised layout-to-image baselines by a large margin.
Tip
Make sure to check out the Stable Diffusion Tips section to learn how to explore the tradeoff between scheduler speed and quality and how to reuse pipeline components efficiently!
If you want to use one of the official checkpoints for a task, explore the gligen Hub organizations!
mindone.diffusers.StableDiffusionGLIGENPipeline
¶
Bases: DiffusionPipeline
, StableDiffusionMixin
Pipeline for text-to-image generation using Stable Diffusion with Grounded-Language-to-Image Generation (GLIGEN).
This model inherits from [DiffusionPipeline
]. Check the superclass documentation for the generic methods the
library implements for all the pipelines (such as downloading or saving, running on a particular device, etc.).
PARAMETER | DESCRIPTION |
---|---|
vae |
Variational Auto-Encoder (VAE) model to encode and decode images to and from latent representations.
TYPE:
|
text_encoder |
Frozen text-encoder (clip-vit-large-patch14).
TYPE:
|
tokenizer |
A
TYPE:
|
unet |
A
TYPE:
|
scheduler |
A scheduler to be used in combination with
TYPE:
|
safety_checker |
Classification module that estimates whether generated images could be considered offensive or harmful. Please refer to the model card for more details about a model's potential harms.
TYPE:
|
feature_extractor |
A
TYPE:
|
Source code in mindone/diffusers/pipelines/stable_diffusion_gligen/pipeline_stable_diffusion_gligen.py
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|
mindone.diffusers.StableDiffusionGLIGENPipeline.__call__(prompt=None, height=None, width=None, num_inference_steps=50, guidance_scale=7.5, gligen_scheduled_sampling_beta=0.3, gligen_phrases=None, gligen_boxes=None, gligen_inpaint_image=None, negative_prompt=None, num_images_per_prompt=1, eta=0.0, generator=None, latents=None, prompt_embeds=None, negative_prompt_embeds=None, output_type='pil', return_dict=False, callback=None, callback_steps=1, cross_attention_kwargs=None, clip_skip=None)
¶
The call function to the pipeline for generation.
PARAMETER | DESCRIPTION |
---|---|
prompt |
The prompt or prompts to guide image generation. If not defined, you need to pass
TYPE:
|
height |
The height in pixels of the generated image.
TYPE:
|
width |
The width in pixels of the generated image.
TYPE:
|
num_inference_steps |
The number of denoising steps. More denoising steps usually lead to a higher quality image at the expense of slower inference.
TYPE:
|
guidance_scale |
A higher guidance scale value encourages the model to generate images closely linked to the text
TYPE:
|
gligen_phrases |
The phrases to guide what to include in each of the regions defined by the corresponding
TYPE:
|
gligen_boxes |
The bounding boxes that identify rectangular regions of the image that are going to be filled with the
content described by the corresponding
TYPE:
|
gligen_inpaint_image |
The input image, if provided, is inpainted with objects described by the
TYPE:
|
gligen_scheduled_sampling_beta |
Scheduled Sampling factor from GLIGEN: Open-Set Grounded Text-to-Image Generation. Scheduled Sampling factor is only varied for scheduled sampling during inference for improved quality and controllability.
TYPE:
|
negative_prompt |
The prompt or prompts to guide what to not include in image generation. If not defined, you need to
pass
TYPE:
|
num_images_per_prompt |
The number of images to generate per prompt.
TYPE:
|
eta |
Corresponds to parameter eta (η) from the DDIM paper. Only applies
to the [
TYPE:
|
generator |
A
TYPE:
|
latents |
Pre-generated noisy latents sampled from a Gaussian distribution, to be used as inputs for image
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents
tensor is generated by sampling using the supplied random
TYPE:
|
prompt_embeds |
Pre-generated text embeddings. Can be used to easily tweak text inputs (prompt weighting). If not
provided, text embeddings are generated from the
TYPE:
|
negative_prompt_embeds |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs (prompt weighting). If
not provided,
TYPE:
|
output_type |
The output format of the generated image. Choose between
TYPE:
|
return_dict |
Whether or not to return a [
TYPE:
|
callback |
A function that calls every
TYPE:
|
callback_steps |
The frequency at which the
TYPE:
|
cross_attention_kwargs |
A kwargs dictionary that if specified is passed along to the [
TYPE:
|
guidance_rescale |
Guidance rescale factor from Common Diffusion Noise Schedules and Sample Steps are Flawed. Guidance rescale factor should fix overexposure when using zero terminal SNR.
TYPE:
|
clip_skip |
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings.
TYPE:
|
RETURNS | DESCRIPTION |
---|---|
[ |
Source code in mindone/diffusers/pipelines/stable_diffusion_gligen/pipeline_stable_diffusion_gligen.py
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|
mindone.diffusers.StableDiffusionGLIGENPipeline.encode_prompt(prompt, num_images_per_prompt, do_classifier_free_guidance, negative_prompt=None, prompt_embeds=None, negative_prompt_embeds=None, lora_scale=None, clip_skip=None)
¶
Encodes the prompt into text encoder hidden states.
PARAMETER | DESCRIPTION |
---|---|
prompt |
prompt to be encoded
TYPE:
|
num_images_per_prompt |
number of images that should be generated per prompt
TYPE:
|
do_classifier_free_guidance |
whether to use classifier free guidance or not
TYPE:
|
negative_prompt |
The prompt or prompts not to guide the image generation. If not defined, one has to pass
TYPE:
|
prompt_embeds |
Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not
provided, text embeddings will be generated from
TYPE:
|
negative_prompt_embeds |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, negative_prompt_embeds will be generated from
TYPE:
|
lora_scale |
A LoRA scale that will be applied to all LoRA layers of the text encoder if LoRA layers are loaded.
TYPE:
|
clip_skip |
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings.
TYPE:
|
Source code in mindone/diffusers/pipelines/stable_diffusion_gligen/pipeline_stable_diffusion_gligen.py
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|
mindone.diffusers.StableDiffusionGLIGENTextImagePipeline
¶
Bases: DiffusionPipeline
, StableDiffusionMixin
Pipeline for text-to-image generation using Stable Diffusion with Grounded-Language-to-Image Generation (GLIGEN).
This model inherits from [DiffusionPipeline
]. Check the superclass documentation for the generic methods the
library implements for all the pipelines (such as downloading or saving, running on a particular device, etc.).
PARAMETER | DESCRIPTION |
---|---|
vae |
Variational Auto-Encoder (VAE) model to encode and decode images to and from latent representations.
TYPE:
|
text_encoder |
Frozen text-encoder (clip-vit-large-patch14).
TYPE:
|
tokenizer |
A
TYPE:
|
processor |
A
TYPE:
|
image_encoder |
Frozen image-encoder (clip-vit-large-patch14).
TYPE:
|
image_project |
A
TYPE:
|
unet |
A
TYPE:
|
scheduler |
A scheduler to be used in combination with
TYPE:
|
safety_checker |
Classification module that estimates whether generated images could be considered offensive or harmful. Please refer to the model card for more details about a model's potential harms.
TYPE:
|
feature_extractor |
A
TYPE:
|
Source code in mindone/diffusers/pipelines/stable_diffusion_gligen/pipeline_stable_diffusion_gligen_text_image.py
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|
mindone.diffusers.StableDiffusionGLIGENTextImagePipeline.__call__(prompt=None, height=None, width=None, num_inference_steps=50, guidance_scale=7.5, gligen_scheduled_sampling_beta=0.3, gligen_phrases=None, gligen_images=None, input_phrases_mask=None, input_images_mask=None, gligen_boxes=None, gligen_inpaint_image=None, negative_prompt=None, num_images_per_prompt=1, eta=0.0, generator=None, latents=None, prompt_embeds=None, negative_prompt_embeds=None, output_type='pil', return_dict=False, callback=None, callback_steps=1, cross_attention_kwargs=None, gligen_normalize_constant=28.7, clip_skip=None)
¶
The call function to the pipeline for generation.
PARAMETER | DESCRIPTION |
---|---|
prompt |
The prompt or prompts to guide image generation. If not defined, you need to pass
TYPE:
|
height |
The height in pixels of the generated image.
TYPE:
|
width |
The width in pixels of the generated image.
TYPE:
|
num_inference_steps |
The number of denoising steps. More denoising steps usually lead to a higher quality image at the expense of slower inference.
TYPE:
|
guidance_scale |
A higher guidance scale value encourages the model to generate images closely linked to the text
TYPE:
|
gligen_phrases |
The phrases to guide what to include in each of the regions defined by the corresponding
TYPE:
|
gligen_images |
The images to guide what to include in each of the regions defined by the corresponding
TYPE:
|
input_phrases_mask |
pre phrases mask input defined by the correspongding
TYPE:
|
input_images_mask |
pre images mask input defined by the correspongding
TYPE:
|
gligen_boxes |
The bounding boxes that identify rectangular regions of the image that are going to be filled with the
content described by the corresponding
TYPE:
|
gligen_inpaint_image |
The input image, if provided, is inpainted with objects described by the
TYPE:
|
gligen_scheduled_sampling_beta |
Scheduled Sampling factor from GLIGEN: Open-Set Grounded Text-to-Image Generation. Scheduled Sampling factor is only varied for scheduled sampling during inference for improved quality and controllability.
TYPE:
|
negative_prompt |
The prompt or prompts to guide what to not include in image generation. If not defined, you need to
pass
TYPE:
|
num_images_per_prompt |
The number of images to generate per prompt.
TYPE:
|
eta |
Corresponds to parameter eta (η) from the DDIM paper. Only applies
to the [
TYPE:
|
generator |
A
TYPE:
|
latents |
Pre-generated noisy latents sampled from a Gaussian distribution, to be used as inputs for image
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents
tensor is generated by sampling using the supplied random
TYPE:
|
prompt_embeds |
Pre-generated text embeddings. Can be used to easily tweak text inputs (prompt weighting). If not
provided, text embeddings are generated from the
TYPE:
|
negative_prompt_embeds |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs (prompt weighting). If
not provided,
TYPE:
|
output_type |
The output format of the generated image. Choose between
TYPE:
|
return_dict |
Whether or not to return a [
TYPE:
|
callback |
A function that calls every
TYPE:
|
callback_steps |
The frequency at which the
TYPE:
|
cross_attention_kwargs |
A kwargs dictionary that if specified is passed along to the [
TYPE:
|
gligen_normalize_constant |
The normalize value of the image embedding.
TYPE:
|
clip_skip |
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings.
TYPE:
|
RETURNS | DESCRIPTION |
---|---|
[ |
Source code in mindone/diffusers/pipelines/stable_diffusion_gligen/pipeline_stable_diffusion_gligen_text_image.py
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|
mindone.diffusers.StableDiffusionGLIGENTextImagePipeline.complete_mask(has_mask, max_objs)
¶
Based on the input mask corresponding value 0 or 1
for each phrases and image, mask the features
corresponding to phrases and images.
Source code in mindone/diffusers/pipelines/stable_diffusion_gligen/pipeline_stable_diffusion_gligen_text_image.py
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|
mindone.diffusers.StableDiffusionGLIGENTextImagePipeline.crop(im, new_width, new_height)
¶
Crop the input image to the specified dimensions.
Source code in mindone/diffusers/pipelines/stable_diffusion_gligen/pipeline_stable_diffusion_gligen_text_image.py
536 537 538 539 540 541 542 543 544 545 |
|
mindone.diffusers.StableDiffusionGLIGENTextImagePipeline.draw_inpaint_mask_from_boxes(boxes, size)
¶
Create an inpainting mask based on given boxes. This function generates an inpainting mask using the provided boxes to mark regions that need to be inpainted.
Source code in mindone/diffusers/pipelines/stable_diffusion_gligen/pipeline_stable_diffusion_gligen_text_image.py
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|
mindone.diffusers.StableDiffusionGLIGENTextImagePipeline.encode_prompt(prompt, num_images_per_prompt, do_classifier_free_guidance, negative_prompt=None, prompt_embeds=None, negative_prompt_embeds=None, lora_scale=None, clip_skip=None)
¶
Encodes the prompt into text encoder hidden states.
PARAMETER | DESCRIPTION |
---|---|
prompt |
prompt to be encoded
TYPE:
|
num_images_per_prompt |
number of images that should be generated per prompt
TYPE:
|
do_classifier_free_guidance |
whether to use classifier free guidance or not
TYPE:
|
negative_prompt |
The prompt or prompts not to guide the image generation. If not defined, one has to pass
TYPE:
|
prompt_embeds |
Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not
provided, text embeddings will be generated from
TYPE:
|
negative_prompt_embeds |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, negative_prompt_embeds will be generated from
TYPE:
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lora_scale |
A LoRA scale that will be applied to all LoRA layers of the text encoder if LoRA layers are loaded.
TYPE:
|
clip_skip |
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings.
TYPE:
|
Source code in mindone/diffusers/pipelines/stable_diffusion_gligen/pipeline_stable_diffusion_gligen_text_image.py
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mindone.diffusers.StableDiffusionGLIGENTextImagePipeline.get_clip_feature(input, normalize_constant, is_image=False)
¶
Get image and phrases embedding by using CLIP pretrain model. The image embedding is transformed into the phrases embedding space through a projection.
Source code in mindone/diffusers/pipelines/stable_diffusion_gligen/pipeline_stable_diffusion_gligen_text_image.py
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mindone.diffusers.StableDiffusionGLIGENTextImagePipeline.get_cross_attention_kwargs_with_grounded(hidden_size, gligen_phrases, gligen_images, gligen_boxes, input_phrases_mask, input_images_mask, repeat_batch, normalize_constant, max_objs)
¶
Prepare the cross-attention kwargs containing information about the grounded input (boxes, mask, image embedding, phrases embedding).
Source code in mindone/diffusers/pipelines/stable_diffusion_gligen/pipeline_stable_diffusion_gligen_text_image.py
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mindone.diffusers.StableDiffusionGLIGENTextImagePipeline.get_cross_attention_kwargs_without_grounded(hidden_size, repeat_batch, max_objs)
¶
Prepare the cross-attention kwargs without information about the grounded input (boxes, mask, image embedding, phrases embedding) (All are zero tensor).
Source code in mindone/diffusers/pipelines/stable_diffusion_gligen/pipeline_stable_diffusion_gligen_text_image.py
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mindone.diffusers.StableDiffusionGLIGENTextImagePipeline.target_size_center_crop(im, new_hw)
¶
Crop and resize the image to the target size while keeping the center.
Source code in mindone/diffusers/pipelines/stable_diffusion_gligen/pipeline_stable_diffusion_gligen_text_image.py
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mindone.diffusers.pipelines.stable_diffusion.StableDiffusionPipelineOutput
dataclass
¶
Bases: BaseOutput
Output class for Stable Diffusion pipelines.
Source code in mindone/diffusers/pipelines/stable_diffusion/pipeline_output.py
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