T2I-Adapter¶
T2I-Adapter: Learning Adapters to Dig out More Controllable Ability for Text-to-Image Diffusion Models by Chong Mou, Xintao Wang, Liangbin Xie, Jian Zhang, Zhongang Qi, Ying Shan, Xiaohu Qie.
Using the pretrained models we can provide control images (for example, a depth map) to control Stable Diffusion text-to-image generation so that it follows the structure of the depth image and fills in the details.
The abstract of the paper is the following:
The incredible generative ability of large-scale text-to-image (T2I) models has demonstrated strong power of learning complex structures and meaningful semantics. However, relying solely on text prompts cannot fully take advantage of the knowledge learned by the model, especially when flexible and accurate controlling (e.g., color and structure) is needed. In this paper, we aim to ``dig out" the capabilities that T2I models have implicitly learned, and then explicitly use them to control the generation more granularly. Specifically, we propose to learn simple and lightweight T2I-Adapters to align internal knowledge in T2I models with external control signals, while freezing the original large T2I models. In this way, we can train various adapters according to different conditions, achieving rich control and editing effects in the color and structure of the generation results. Further, the proposed T2I-Adapters have attractive properties of practical value, such as composability and generalization ability. Extensive experiments demonstrate that our T2I-Adapter has promising generation quality and a wide range of applications.
mindone.diffusers.StableDiffusionAdapterPipeline
¶
Bases: DiffusionPipeline
, StableDiffusionMixin
Pipeline for text-to-image generation using Stable Diffusion augmented with T2I-Adapter https://arxiv.org/abs/2302.08453
This model inherits from [DiffusionPipeline
]. Check the superclass documentation for the generic methods the
library implements for all the pipelines (such as downloading or saving, running on a particular device, etc.)
PARAMETER | DESCRIPTION |
---|---|
adapter |
Provides additional conditioning to the unet during the denoising process. If you set multiple Adapter as a list, the outputs from each Adapter are added together to create one combined additional conditioning.
TYPE:
|
adapter_weights |
List of floats representing the weight which will be multiply to each adapter's output before adding them together.
TYPE:
|
vae |
Variational Auto-Encoder (VAE) Model to encode and decode images to and from latent representations.
TYPE:
|
text_encoder |
Frozen text-encoder. Stable Diffusion uses the text portion of CLIP, specifically the clip-vit-large-patch14 variant.
TYPE:
|
tokenizer |
Tokenizer of class CLIPTokenizer.
TYPE:
|
unet |
Conditional U-Net architecture to denoise the encoded image latents.
TYPE:
|
scheduler |
A scheduler to be used in combination with
TYPE:
|
safety_checker |
Classification module that estimates whether generated images could be considered offensive or harmful. Please, refer to the model card for details.
TYPE:
|
feature_extractor |
Model that extracts features from generated images to be used as inputs for the
TYPE:
|
Source code in mindone/diffusers/pipelines/t2i_adapter/pipeline_stable_diffusion_adapter.py
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|
mindone.diffusers.StableDiffusionAdapterPipeline.__call__(prompt=None, image=None, height=None, width=None, num_inference_steps=50, timesteps=None, sigmas=None, guidance_scale=7.5, negative_prompt=None, num_images_per_prompt=1, eta=0.0, generator=None, latents=None, prompt_embeds=None, negative_prompt_embeds=None, output_type='pil', return_dict=False, callback=None, callback_steps=1, cross_attention_kwargs=None, adapter_conditioning_scale=1.0, clip_skip=None)
¶
Function invoked when calling the pipeline for generation.
PARAMETER | DESCRIPTION |
---|---|
prompt |
The prompt or prompts to guide the image generation. If not defined, one has to pass
TYPE:
|
image |
The Adapter input condition. Adapter uses this input condition to generate guidance to Unet. If the
type is specified as
TYPE:
|
height |
The height in pixels of the generated image.
TYPE:
|
width |
The width in pixels of the generated image.
TYPE:
|
num_inference_steps |
The number of denoising steps. More denoising steps usually lead to a higher quality image at the expense of slower inference.
TYPE:
|
timesteps |
Custom timesteps to use for the denoising process with schedulers which support a
TYPE:
|
sigmas |
Custom sigmas to use for the denoising process with schedulers which support a
TYPE:
|
guidance_scale |
Guidance scale as defined in Classifier-Free Diffusion Guidance.
TYPE:
|
negative_prompt |
The prompt or prompts not to guide the image generation. If not defined, one has to pass
TYPE:
|
num_images_per_prompt |
The number of images to generate per prompt.
TYPE:
|
eta |
Corresponds to parameter eta (η) in the DDIM paper: https://arxiv.org/abs/2010.02502. Only applies to
[
TYPE:
|
generator |
One or a list of np.random.Generator to make generation deterministic.
TYPE:
|
latents |
Pre-generated noisy latents, sampled from a Gaussian distribution, to be used as inputs for image
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents
tensor will ge generated by sampling using the supplied random
TYPE:
|
prompt_embeds |
Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not
provided, text embeddings will be generated from
TYPE:
|
negative_prompt_embeds |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, negative_prompt_embeds will be generated from
TYPE:
|
output_type |
The output format of the generate image. Choose between
PIL:
TYPE:
|
return_dict |
Whether or not to return a [
TYPE:
|
callback |
A function that will be called every
TYPE:
|
callback_steps |
The frequency at which the
TYPE:
|
cross_attention_kwargs |
A kwargs dictionary that if specified is passed along to the
TYPE:
|
adapter_conditioning_scale |
The outputs of the adapter are multiplied by
TYPE:
|
clip_skip |
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings.
TYPE:
|
RETURNS | DESCRIPTION |
---|---|
[ |
|
[ |
|
`tuple. When returning a tuple, the first element is a list with the generated images, and the second |
|
element is a list of |
|
"not-safe-for-work" (nsfw) content, according to the |
Source code in mindone/diffusers/pipelines/t2i_adapter/pipeline_stable_diffusion_adapter.py
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|
mindone.diffusers.StableDiffusionAdapterPipeline.encode_prompt(prompt, num_images_per_prompt, do_classifier_free_guidance, negative_prompt=None, prompt_embeds=None, negative_prompt_embeds=None, lora_scale=None, clip_skip=None)
¶
Encodes the prompt into text encoder hidden states.
PARAMETER | DESCRIPTION |
---|---|
prompt |
prompt to be encoded
TYPE:
|
num_images_per_prompt |
number of images that should be generated per prompt
TYPE:
|
do_classifier_free_guidance |
whether to use classifier free guidance or not
TYPE:
|
negative_prompt |
The prompt or prompts not to guide the image generation. If not defined, one has to pass
TYPE:
|
prompt_embeds |
Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not
provided, text embeddings will be generated from
TYPE:
|
negative_prompt_embeds |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, negative_prompt_embeds will be generated from
TYPE:
|
lora_scale |
A LoRA scale that will be applied to all LoRA layers of the text encoder if LoRA layers are loaded.
TYPE:
|
clip_skip |
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings.
TYPE:
|
Source code in mindone/diffusers/pipelines/t2i_adapter/pipeline_stable_diffusion_adapter.py
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|
mindone.diffusers.StableDiffusionAdapterPipeline.get_guidance_scale_embedding(w, embedding_dim=512, dtype=ms.float32)
¶
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
PARAMETER | DESCRIPTION |
---|---|
w |
Generate embedding vectors with a specified guidance scale to subsequently enrich timestep embeddings.
TYPE:
|
embedding_dim |
Dimension of the embeddings to generate.
TYPE:
|
dtype |
Data type of the generated embeddings.
TYPE:
|
RETURNS | DESCRIPTION |
---|---|
Tensor
|
|
Source code in mindone/diffusers/pipelines/t2i_adapter/pipeline_stable_diffusion_adapter.py
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|
mindone.diffusers.StableDiffusionXLAdapterPipeline
¶
Bases: DiffusionPipeline
, StableDiffusionMixin
, TextualInversionLoaderMixin
, StableDiffusionXLLoraLoaderMixin
, IPAdapterMixin
, FromSingleFileMixin
Pipeline for text-to-image generation using Stable Diffusion augmented with T2I-Adapter https://arxiv.org/abs/2302.08453
This model inherits from [DiffusionPipeline
]. Check the superclass documentation for the generic methods the
library implements for all the pipelines (such as downloading or saving, running on a particular device, etc.)
The pipeline also inherits the following loading methods
- [
~loaders.TextualInversionLoaderMixin.load_textual_inversion
] for loading textual inversion embeddings - [
~loaders.FromSingleFileMixin.from_single_file
] for loading.ckpt
files - [
~loaders.StableDiffusionXLLoraLoaderMixin.load_lora_weights
] for loading LoRA weights - [
~loaders.StableDiffusionXLLoraLoaderMixin.save_lora_weights
] for saving LoRA weights - [
~loaders.IPAdapterMixin.load_ip_adapter
] for loading IP Adapters
PARAMETER | DESCRIPTION |
---|---|
adapter |
Provides additional conditioning to the unet during the denoising process. If you set multiple Adapter as a list, the outputs from each Adapter are added together to create one combined additional conditioning.
TYPE:
|
adapter_weights |
List of floats representing the weight which will be multiply to each adapter's output before adding them together.
TYPE:
|
vae |
Variational Auto-Encoder (VAE) Model to encode and decode images to and from latent representations.
TYPE:
|
text_encoder |
Frozen text-encoder. Stable Diffusion uses the text portion of CLIP, specifically the clip-vit-large-patch14 variant.
TYPE:
|
tokenizer |
Tokenizer of class CLIPTokenizer.
TYPE:
|
unet |
Conditional U-Net architecture to denoise the encoded image latents.
TYPE:
|
scheduler |
A scheduler to be used in combination with
TYPE:
|
safety_checker |
Classification module that estimates whether generated images could be considered offensive or harmful. Please, refer to the model card for details.
TYPE:
|
feature_extractor |
Model that extracts features from generated images to be used as inputs for the
TYPE:
|
Source code in mindone/diffusers/pipelines/t2i_adapter/pipeline_stable_diffusion_xl_adapter.py
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|
mindone.diffusers.StableDiffusionXLAdapterPipeline.__call__(prompt=None, prompt_2=None, image=None, height=None, width=None, num_inference_steps=50, sigmas=None, timesteps=None, denoising_end=None, guidance_scale=5.0, negative_prompt=None, negative_prompt_2=None, num_images_per_prompt=1, eta=0.0, generator=None, latents=None, prompt_embeds=None, negative_prompt_embeds=None, pooled_prompt_embeds=None, negative_pooled_prompt_embeds=None, ip_adapter_image=None, ip_adapter_image_embeds=None, output_type='pil', return_dict=False, callback=None, callback_steps=1, cross_attention_kwargs=None, guidance_rescale=0.0, original_size=None, crops_coords_top_left=(0, 0), target_size=None, negative_original_size=None, negative_crops_coords_top_left=(0, 0), negative_target_size=None, adapter_conditioning_scale=1.0, adapter_conditioning_factor=1.0, clip_skip=None)
¶
Function invoked when calling the pipeline for generation.
PARAMETER | DESCRIPTION |
---|---|
prompt |
The prompt or prompts to guide the image generation. If not defined, one has to pass
TYPE:
|
prompt_2 |
The prompt or prompts to be sent to the
TYPE:
|
image |
The Adapter input condition. Adapter uses this input condition to generate guidance to Unet. If the
type is specified as
TYPE:
|
height |
The height in pixels of the generated image. Anything below 512 pixels won't work well for stabilityai/stable-diffusion-xl-base-1.0 and checkpoints that are not specifically fine-tuned on low resolutions.
TYPE:
|
width |
The width in pixels of the generated image. Anything below 512 pixels won't work well for stabilityai/stable-diffusion-xl-base-1.0 and checkpoints that are not specifically fine-tuned on low resolutions.
TYPE:
|
num_inference_steps |
The number of denoising steps. More denoising steps usually lead to a higher quality image at the expense of slower inference.
TYPE:
|
timesteps |
Custom timesteps to use for the denoising process with schedulers which support a
TYPE:
|
sigmas |
Custom sigmas to use for the denoising process with schedulers which support a
TYPE:
|
denoising_end |
When specified, determines the fraction (between 0.0 and 1.0) of the total denoising process to be completed before it is intentionally prematurely terminated. As a result, the returned sample will still retain a substantial amount of noise as determined by the discrete timesteps selected by the scheduler. The denoising_end parameter should ideally be utilized when this pipeline forms a part of a "Mixture of Denoisers" multi-pipeline setup, as elaborated in Refining the Image Output
TYPE:
|
guidance_scale |
Guidance scale as defined in Classifier-Free Diffusion Guidance.
TYPE:
|
negative_prompt |
The prompt or prompts not to guide the image generation. If not defined, one has to pass
TYPE:
|
negative_prompt_2 |
The prompt or prompts not to guide the image generation to be sent to
TYPE:
|
num_images_per_prompt |
The number of images to generate per prompt.
TYPE:
|
eta |
Corresponds to parameter eta (η) in the DDIM paper: https://arxiv.org/abs/2010.02502. Only applies to
[
TYPE:
|
generator |
One or a list of numpy generator(s) to make generation deterministic.
TYPE:
|
latents |
Pre-generated noisy latents, sampled from a Gaussian distribution, to be used as inputs for image
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents
tensor will ge generated by sampling using the supplied random
TYPE:
|
prompt_embeds |
Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not
provided, text embeddings will be generated from
TYPE:
|
negative_prompt_embeds |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, negative_prompt_embeds will be generated from
TYPE:
|
pooled_prompt_embeds |
Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting.
If not provided, pooled text embeddings will be generated from
TYPE:
|
negative_pooled_prompt_embeds |
Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, pooled negative_prompt_embeds will be generated from
TYPE:
|
ip_adapter_image |
(
TYPE:
|
ip_adapter_image_embeds |
Pre-generated image embeddings for IP-Adapter. It should be a list of length same as number of IP-adapters.
Each element should be a tensor of shape
TYPE:
|
output_type |
The output format of the generate image. Choose between
PIL:
TYPE:
|
return_dict |
Whether or not to return a [
TYPE:
|
callback |
A function that will be called every
TYPE:
|
callback_steps |
The frequency at which the
TYPE:
|
cross_attention_kwargs |
A kwargs dictionary that if specified is passed along to the
TYPE:
|
guidance_rescale |
Guidance rescale factor proposed by Common Diffusion Noise Schedules and Sample Steps are
Flawed
TYPE:
|
original_size |
If
TYPE:
|
crops_coords_top_left |
TYPE:
|
target_size |
For most cases,
TYPE:
|
negative_original_size |
To negatively condition the generation process based on a specific image resolution. Part of SDXL's micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. For more information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208.
TYPE:
|
negative_crops_coords_top_left |
To negatively condition the generation process based on a specific crop coordinates. Part of SDXL's micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. For more information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208.
TYPE:
|
negative_target_size |
To negatively condition the generation process based on a target image resolution. It should be as same
as the
TYPE:
|
adapter_conditioning_scale |
The outputs of the adapter are multiplied by
TYPE:
|
adapter_conditioning_factor |
The fraction of timesteps for which adapter should be applied. If
TYPE:
|
clip_skip |
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings.
TYPE:
|
RETURNS | DESCRIPTION |
---|---|
[ |
|
[ |
|
|
Source code in mindone/diffusers/pipelines/t2i_adapter/pipeline_stable_diffusion_xl_adapter.py
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mindone.diffusers.StableDiffusionXLAdapterPipeline.encode_prompt(prompt, prompt_2=None, num_images_per_prompt=1, do_classifier_free_guidance=True, negative_prompt=None, negative_prompt_2=None, prompt_embeds=None, negative_prompt_embeds=None, pooled_prompt_embeds=None, negative_pooled_prompt_embeds=None, lora_scale=None, clip_skip=None)
¶
Encodes the prompt into text encoder hidden states.
PARAMETER | DESCRIPTION |
---|---|
prompt |
prompt to be encoded
TYPE:
|
prompt_2 |
The prompt or prompts to be sent to the
TYPE:
|
num_images_per_prompt |
number of images that should be generated per prompt
TYPE:
|
do_classifier_free_guidance |
whether to use classifier free guidance or not
TYPE:
|
negative_prompt |
The prompt or prompts not to guide the image generation. If not defined, one has to pass
TYPE:
|
negative_prompt_2 |
The prompt or prompts not to guide the image generation to be sent to
TYPE:
|
prompt_embeds |
Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not
provided, text embeddings will be generated from
TYPE:
|
negative_prompt_embeds |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, negative_prompt_embeds will be generated from
TYPE:
|
pooled_prompt_embeds |
Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting.
If not provided, pooled text embeddings will be generated from
TYPE:
|
negative_pooled_prompt_embeds |
Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, pooled negative_prompt_embeds will be generated from
TYPE:
|
lora_scale |
A lora scale that will be applied to all LoRA layers of the text encoder if LoRA layers are loaded.
TYPE:
|
clip_skip |
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings.
TYPE:
|
Source code in mindone/diffusers/pipelines/t2i_adapter/pipeline_stable_diffusion_xl_adapter.py
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|
mindone.diffusers.StableDiffusionXLAdapterPipeline.get_guidance_scale_embedding(w, embedding_dim=512, dtype=ms.float32)
¶
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
PARAMETER | DESCRIPTION |
---|---|
w |
Generate embedding vectors with a specified guidance scale to subsequently enrich timestep embeddings.
TYPE:
|
embedding_dim |
Dimension of the embeddings to generate.
TYPE:
|
dtype |
Data type of the generated embeddings.
TYPE:
|
RETURNS | DESCRIPTION |
---|---|
Tensor
|
|
Source code in mindone/diffusers/pipelines/t2i_adapter/pipeline_stable_diffusion_xl_adapter.py
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