InstructPix2Pix¶
InstructPix2Pix: Learning to Follow Image Editing Instructions is by Tim Brooks, Aleksander Holynski and Alexei A. Efros.
The abstract from the paper is:
We propose a method for editing images from human instructions: given an input image and a written instruction that tells the model what to do, our model follows these instructions to edit the image. To obtain training data for this problem, we combine the knowledge of two large pretrained models -- a language model (GPT-3) and a text-to-image model (Stable Diffusion) -- to generate a large dataset of image editing examples. Our conditional diffusion model, InstructPix2Pix, is trained on our generated data, and generalizes to real images and user-written instructions at inference time. Since it performs edits in the forward pass and does not require per example fine-tuning or inversion, our model edits images quickly, in a matter of seconds. We show compelling editing results for a diverse collection of input images and written instructions.
You can find additional information about InstructPix2Pix on the project page, original codebase, and try it out in a demo.
Tip
Make sure to check out the Schedulers guide to learn how to explore the tradeoff between scheduler speed and quality, and see the reuse components across pipelines section to learn how to efficiently load the same components into multiple pipelines.
mindone.diffusers.StableDiffusionInstructPix2PixPipeline
¶
Bases: DiffusionPipeline
, StableDiffusionMixin
, TextualInversionLoaderMixin
, LoraLoaderMixin
, IPAdapterMixin
Pipeline for pixel-level image editing by following text instructions (based on Stable Diffusion).
This model inherits from [DiffusionPipeline
]. Check the superclass documentation for the generic methods
implemented for all pipelines (downloading, saving, running on a particular device, etc.).
The pipeline also inherits the following loading methods
- [
~loaders.TextualInversionLoaderMixin.load_textual_inversion
] for loading textual inversion embeddings - [
~loaders.LoraLoaderMixin.load_lora_weights
] for loading LoRA weights - [
~loaders.LoraLoaderMixin.save_lora_weights
] for saving LoRA weights - [
~loaders.IPAdapterMixin.load_ip_adapter
] for loading IP Adapters
PARAMETER | DESCRIPTION |
---|---|
vae |
Variational Auto-Encoder (VAE) model to encode and decode images to and from latent representations.
TYPE:
|
text_encoder |
Frozen text-encoder (clip-vit-large-patch14).
TYPE:
|
tokenizer |
A
TYPE:
|
unet |
A
TYPE:
|
scheduler |
A scheduler to be used in combination with
TYPE:
|
safety_checker |
Classification module that estimates whether generated images could be considered offensive or harmful. Please refer to the model card for more details about a model's potential harms.
TYPE:
|
feature_extractor |
A
TYPE:
|
Source code in mindone/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_instruct_pix2pix.py
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|
mindone.diffusers.StableDiffusionInstructPix2PixPipeline.__call__(prompt=None, image=None, num_inference_steps=100, guidance_scale=7.5, image_guidance_scale=1.5, negative_prompt=None, num_images_per_prompt=1, eta=0.0, generator=None, latents=None, prompt_embeds=None, negative_prompt_embeds=None, ip_adapter_image=None, ip_adapter_image_embeds=None, output_type='pil', return_dict=False, callback_on_step_end=None, callback_on_step_end_tensor_inputs=['latents'], cross_attention_kwargs=None, **kwargs)
¶
The call function to the pipeline for generation.
PARAMETER | DESCRIPTION |
---|---|
prompt |
The prompt or prompts to guide image generation. If not defined, you need to pass
TYPE:
|
image |
TYPE:
|
num_inference_steps |
The number of denoising steps. More denoising steps usually lead to a higher quality image at the expense of slower inference.
TYPE:
|
guidance_scale |
A higher guidance scale value encourages the model to generate images closely linked to the text
TYPE:
|
image_guidance_scale |
Push the generated image towards the initial
TYPE:
|
negative_prompt |
The prompt or prompts to guide what to not include in image generation. If not defined, you need to
pass
TYPE:
|
num_images_per_prompt |
The number of images to generate per prompt.
TYPE:
|
eta |
Corresponds to parameter eta (η) from the DDIM paper. Only applies
to the [
TYPE:
|
generator |
A
TYPE:
|
latents |
Pre-generated noisy latents sampled from a Gaussian distribution, to be used as inputs for image
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents
tensor is generated by sampling using the supplied random
TYPE:
|
prompt_embeds |
Pre-generated text embeddings. Can be used to easily tweak text inputs (prompt weighting). If not
provided, text embeddings are generated from the
TYPE:
|
negative_prompt_embeds |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs (prompt weighting). If
not provided,
TYPE:
|
ip_adapter_image |
(
TYPE:
|
output_type |
The output format of the generated image. Choose between
TYPE:
|
return_dict |
Whether or not to return a [
TYPE:
|
callback_on_step_end |
A function or a subclass of
TYPE:
|
callback_on_step_end_tensor_inputs |
The list of tensor inputs for the
TYPE:
|
cross_attention_kwargs |
A kwargs dictionary that if specified is passed along to the [
TYPE:
|
>>> import PIL
>>> import requests
>>> import mindspore as ms
>>> from io import BytesIO
>>> from mindone.diffusers import StableDiffusionInstructPix2PixPipeline
>>> def download_image(url):
... response = requests.get(url)
... return PIL.Image.open(BytesIO(response.content)).convert("RGB")
>>> img_url = "https://huggingface.co/datasets/diffusers/diffusers-images-docs/resolve/main/mountain.png"
>>> image = download_image(img_url).resize((512, 512))
>>> pipe = StableDiffusionInstructPix2PixPipeline.from_pretrained(
... "timbrooks/instruct-pix2pix", mindspore_dtype=ms.float16
... )
>>> prompt = "make the mountains snowy"
>>> image = pipe(prompt=prompt, image=image)[0][0]
RETURNS | DESCRIPTION |
---|---|
[ |
Source code in mindone/diffusers/pipelines/stable_diffusion/pipeline_stable_diffusion_instruct_pix2pix.py
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|
mindone.diffusers.StableDiffusionXLInstructPix2PixPipeline
¶
Bases: DiffusionPipeline
, StableDiffusionMixin
, TextualInversionLoaderMixin
, FromSingleFileMixin
, StableDiffusionXLLoraLoaderMixin
Pipeline for pixel-level image editing by following text instructions. Based on Stable Diffusion XL.
This model inherits from [DiffusionPipeline
]. Check the superclass documentation for the generic methods the
library implements for all the pipelines (such as downloading or saving, running on a particular device, etc.)
The pipeline also inherits the following loading methods
- [
~loaders.TextualInversionLoaderMixin.load_textual_inversion
] for loading textual inversion embeddings - [
~loaders.FromSingleFileMixin.from_single_file
] for loading.ckpt
files - [
~loaders.StableDiffusionXLLoraLoaderMixin.load_lora_weights
] for loading LoRA weights - [
~loaders.StableDiffusionXLLoraLoaderMixin.save_lora_weights
] for saving LoRA weights
PARAMETER | DESCRIPTION |
---|---|
vae |
Variational Auto-Encoder (VAE) Model to encode and decode images to and from latent representations.
TYPE:
|
text_encoder |
Frozen text-encoder. Stable Diffusion XL uses the text portion of CLIP, specifically the clip-vit-large-patch14 variant.
TYPE:
|
text_encoder_2 |
Second frozen text-encoder. Stable Diffusion XL uses the text and pool portion of CLIP, specifically the laion/CLIP-ViT-bigG-14-laion2B-39B-b160k variant.
TYPE:
|
tokenizer |
Tokenizer of class CLIPTokenizer.
TYPE:
|
tokenizer_2 |
Second Tokenizer of class CLIPTokenizer.
TYPE:
|
unet |
Conditional U-Net architecture to denoise the encoded image latents.
TYPE:
|
scheduler |
A scheduler to be used in combination with
TYPE:
|
requires_aesthetics_score |
Whether the
TYPE:
|
force_zeros_for_empty_prompt |
Whether the negative prompt embeddings shall be forced to always be set to 0. Also see the config of
TYPE:
|
add_watermarker |
Whether to use the invisible_watermark library to watermark output images. If not defined, it will default to True if the package is installed, otherwise no watermarker will be used.
TYPE:
|
is_cosxl_edit |
When set the image latents are scaled.
TYPE:
|
Source code in mindone/diffusers/pipelines/stable_diffusion_xl/pipeline_stable_diffusion_xl_instruct_pix2pix.py
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|
mindone.diffusers.StableDiffusionXLInstructPix2PixPipeline.__call__(prompt=None, prompt_2=None, image=None, height=None, width=None, num_inference_steps=100, denoising_end=None, guidance_scale=5.0, image_guidance_scale=1.5, negative_prompt=None, negative_prompt_2=None, num_images_per_prompt=1, eta=0.0, generator=None, latents=None, prompt_embeds=None, negative_prompt_embeds=None, pooled_prompt_embeds=None, negative_pooled_prompt_embeds=None, output_type='pil', return_dict=False, callback=None, callback_steps=1, cross_attention_kwargs=None, guidance_rescale=0.0, original_size=None, crops_coords_top_left=(0, 0), target_size=None)
¶
Function invoked when calling the pipeline for generation.
PARAMETER | DESCRIPTION |
---|---|
prompt |
The prompt or prompts to guide the image generation. If not defined, one has to pass
TYPE:
|
prompt_2 |
The prompt or prompts to be sent to the
TYPE:
|
image |
The image(s) to modify with the pipeline.
TYPE:
|
height |
The height in pixels of the generated image.
TYPE:
|
width |
The width in pixels of the generated image.
TYPE:
|
num_inference_steps |
The number of denoising steps. More denoising steps usually lead to a higher quality image at the expense of slower inference.
TYPE:
|
denoising_end |
When specified, determines the fraction (between 0.0 and 1.0) of the total denoising process to be completed before it is intentionally prematurely terminated. As a result, the returned sample will still retain a substantial amount of noise as determined by the discrete timesteps selected by the scheduler. The denoising_end parameter should ideally be utilized when this pipeline forms a part of a "Mixture of Denoisers" multi-pipeline setup, as elaborated in Refining the Image Output
TYPE:
|
guidance_scale |
Guidance scale as defined in Classifier-Free Diffusion Guidance.
TYPE:
|
image_guidance_scale |
Image guidance scale is to push the generated image towards the inital image
TYPE:
|
negative_prompt |
The prompt or prompts not to guide the image generation. If not defined, one has to pass
TYPE:
|
negative_prompt_2 |
The prompt or prompts not to guide the image generation to be sent to
TYPE:
|
num_images_per_prompt |
The number of images to generate per prompt.
TYPE:
|
eta |
Corresponds to parameter eta (η) in the DDIM paper: https://arxiv.org/abs/2010.02502. Only applies to
[
TYPE:
|
generator |
One or a list of numpy generator(s) to make generation deterministic.
TYPE:
|
latents |
Pre-generated noisy latents, sampled from a Gaussian distribution, to be used as inputs for image
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents
tensor will ge generated by sampling using the supplied random
TYPE:
|
prompt_embeds |
Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not
provided, text embeddings will be generated from
TYPE:
|
negative_prompt_embeds |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, negative_prompt_embeds will be generated from
TYPE:
|
pooled_prompt_embeds |
Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting.
If not provided, pooled text embeddings will be generated from
TYPE:
|
negative_pooled_prompt_embeds |
Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, pooled negative_prompt_embeds will be generated from
TYPE:
|
output_type |
The output format of the generate image. Choose between
PIL:
TYPE:
|
return_dict |
Whether or not to return a [
TYPE:
|
callback |
A function that will be called every
TYPE:
|
callback_steps |
The frequency at which the
TYPE:
|
cross_attention_kwargs |
A kwargs dictionary that if specified is passed along to the
TYPE:
|
guidance_rescale |
Guidance rescale factor proposed by Common Diffusion Noise Schedules and Sample Steps are
Flawed
TYPE:
|
original_size |
If
TYPE:
|
crops_coords_top_left |
TYPE:
|
target_size |
For most cases,
TYPE:
|
aesthetic_score |
Used to simulate an aesthetic score of the generated image by influencing the positive text condition. Part of SDXL's micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952.
TYPE:
|
negative_aesthetic_score |
Part of SDXL's micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. Can be used to simulate an aesthetic score of the generated image by influencing the negative text condition.
TYPE:
|
RETURNS | DESCRIPTION |
---|---|
[ |
|
[ |
|
|
Source code in mindone/diffusers/pipelines/stable_diffusion_xl/pipeline_stable_diffusion_xl_instruct_pix2pix.py
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|
mindone.diffusers.StableDiffusionXLInstructPix2PixPipeline.encode_prompt(prompt, prompt_2=None, num_images_per_prompt=1, do_classifier_free_guidance=True, negative_prompt=None, negative_prompt_2=None, prompt_embeds=None, negative_prompt_embeds=None, pooled_prompt_embeds=None, negative_pooled_prompt_embeds=None, lora_scale=None)
¶
Encodes the prompt into text encoder hidden states.
PARAMETER | DESCRIPTION |
---|---|
prompt |
prompt to be encoded
TYPE:
|
prompt_2 |
The prompt or prompts to be sent to the
TYPE:
|
num_images_per_prompt |
number of images that should be generated per prompt
TYPE:
|
do_classifier_free_guidance |
whether to use classifier free guidance or not
TYPE:
|
negative_prompt |
The prompt or prompts not to guide the image generation. If not defined, one has to pass
TYPE:
|
negative_prompt_2 |
The prompt or prompts not to guide the image generation to be sent to
TYPE:
|
prompt_embeds |
Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not
provided, text embeddings will be generated from
TYPE:
|
negative_prompt_embeds |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, negative_prompt_embeds will be generated from
TYPE:
|
pooled_prompt_embeds |
Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting.
If not provided, pooled text embeddings will be generated from
TYPE:
|
negative_pooled_prompt_embeds |
Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, pooled negative_prompt_embeds will be generated from
TYPE:
|
lora_scale |
A lora scale that will be applied to all LoRA layers of the text encoder if LoRA layers are loaded.
TYPE:
|
Source code in mindone/diffusers/pipelines/stable_diffusion_xl/pipeline_stable_diffusion_xl_instruct_pix2pix.py
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|