ControlNet-XS with Stable Diffusion XL¶
ControlNet-XS was introduced in ControlNet-XS by Denis Zavadski and Carsten Rother. It is based on the observation that the control model in the original ControlNet can be made much smaller and still produce good results.
Like the original ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.
ControlNet-XS generates images with comparable quality to a regular ControlNet, but it is 20-25% faster (see benchmark) and uses ~45% less memory.
Here's the overview from the project page:
With increasing computing capabilities, current model architectures appear to follow the trend of simply upscaling all components without validating the necessity for doing so. In this project we investigate the size and architectural design of ControlNet [Zhang et al., 2023] for controlling the image generation process with stable diffusion-based models. We show that a new architecture with as little as 1% of the parameters of the base model achieves state-of-the art results, considerably better than ControlNet in terms of FID score. Hence we call it ControlNet-XS. We provide the code for controlling StableDiffusion-XL [Podell et al., 2023] (Model B, 48M Parameters) and StableDiffusion 2.1 [Rombach et al. 2022] (Model B, 14M Parameters), all under openrail license.
This model was contributed by UmerHA. ❤️
Tip
🧪 Many of the SDXL ControlNet checkpoints are experimental, and there is a lot of room for improvement. Feel free to open an Issue and leave us feedback on how we can improve!
Tip
Make sure to check out the Schedulers guide to learn how to explore the tradeoff between scheduler speed and quality, and see the reuse components across pipelines section to learn how to efficiently load the same components into multiple pipelines.
mindone.diffusers.StableDiffusionXLControlNetXSPipeline
¶
Bases: DiffusionPipeline
, TextualInversionLoaderMixin
, StableDiffusionXLLoraLoaderMixin
, FromSingleFileMixin
Pipeline for text-to-image generation using Stable Diffusion XL with ControlNet-XS guidance.
This model inherits from [DiffusionPipeline
]. Check the superclass documentation for the generic methods
implemented for all pipelines (downloading, saving, running on a particular device, etc.).
The pipeline also inherits the following loading methods
- [
~loaders.TextualInversionLoaderMixin.load_textual_inversion
] for loading textual inversion embeddings - [
loaders.StableDiffusionXLLoraLoaderMixin.load_lora_weights
] for loading LoRA weights - [
loaders.FromSingleFileMixin.from_single_file
] for loading.ckpt
files
PARAMETER | DESCRIPTION |
---|---|
vae |
Variational Auto-Encoder (VAE) model to encode and decode images to and from latent representations.
TYPE:
|
text_encoder |
Frozen text-encoder (clip-vit-large-patch14).
TYPE:
|
text_encoder_2 |
Second frozen text-encoder (laion/CLIP-ViT-bigG-14-laion2B-39B-b160k).
TYPE:
|
tokenizer |
A
TYPE:
|
tokenizer_2 |
A
TYPE:
|
unet |
A [
TYPE:
|
controlnet |
A [
TYPE:
|
scheduler |
A scheduler to be used in combination with
TYPE:
|
force_zeros_for_empty_prompt |
Whether the negative prompt embeddings should always be set to 0. Also see the config of
TYPE:
|
add_watermarker |
Whether to use the invisible_watermark library to
watermark output images. If not defined, it defaults to
TYPE:
|
Source code in mindone/diffusers/pipelines/controlnet_xs/pipeline_controlnet_xs_sd_xl.py
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|
mindone.diffusers.StableDiffusionXLControlNetXSPipeline.__call__(prompt=None, prompt_2=None, image=None, height=None, width=None, num_inference_steps=50, guidance_scale=5.0, negative_prompt=None, negative_prompt_2=None, num_images_per_prompt=1, eta=0.0, generator=None, latents=None, prompt_embeds=None, negative_prompt_embeds=None, pooled_prompt_embeds=None, negative_pooled_prompt_embeds=None, output_type='pil', return_dict=False, cross_attention_kwargs=None, controlnet_conditioning_scale=1.0, control_guidance_start=0.0, control_guidance_end=1.0, original_size=None, crops_coords_top_left=(0, 0), target_size=None, negative_original_size=None, negative_crops_coords_top_left=(0, 0), negative_target_size=None, clip_skip=None, callback_on_step_end=None, callback_on_step_end_tensor_inputs=['latents'])
¶
The call function to the pipeline for generation.
PARAMETER | DESCRIPTION |
---|---|
prompt |
The prompt or prompts to guide image generation. If not defined, you need to pass
TYPE:
|
prompt_2 |
The prompt or prompts to be sent to
TYPE:
|
image |
The ControlNet input condition to provide guidance to the
TYPE:
|
height |
The height in pixels of the generated image. Anything below 512 pixels won't work well for stabilityai/stable-diffusion-xl-base-1.0 and checkpoints that are not specifically fine-tuned on low resolutions.
TYPE:
|
width |
The width in pixels of the generated image. Anything below 512 pixels won't work well for stabilityai/stable-diffusion-xl-base-1.0 and checkpoints that are not specifically fine-tuned on low resolutions.
TYPE:
|
num_inference_steps |
The number of denoising steps. More denoising steps usually lead to a higher quality image at the expense of slower inference.
TYPE:
|
guidance_scale |
A higher guidance scale value encourages the model to generate images closely linked to the text
TYPE:
|
negative_prompt |
The prompt or prompts to guide what to not include in image generation. If not defined, you need to
pass
TYPE:
|
negative_prompt_2 |
The prompt or prompts to guide what to not include in image generation. This is sent to
TYPE:
|
num_images_per_prompt |
The number of images to generate per prompt.
TYPE:
|
eta |
Corresponds to parameter eta (η) from the DDIM paper. Only applies
to the [
TYPE:
|
generator |
A
TYPE:
|
latents |
Pre-generated noisy latents sampled from a Gaussian distribution, to be used as inputs for image
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents
tensor is generated by sampling using the supplied random
TYPE:
|
prompt_embeds |
Pre-generated text embeddings. Can be used to easily tweak text inputs (prompt weighting). If not
provided, text embeddings are generated from the
TYPE:
|
negative_prompt_embeds |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs (prompt weighting). If
not provided,
TYPE:
|
pooled_prompt_embeds |
Pre-generated pooled text embeddings. Can be used to easily tweak text inputs (prompt weighting). If
not provided, pooled text embeddings are generated from
TYPE:
|
negative_pooled_prompt_embeds |
Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs (prompt
weighting). If not provided, pooled
TYPE:
|
output_type |
The output format of the generated image. Choose between
TYPE:
|
return_dict |
Whether or not to return a [
TYPE:
|
cross_attention_kwargs |
A kwargs dictionary that if specified is passed along to the [
TYPE:
|
controlnet_conditioning_scale |
The outputs of the ControlNet are multiplied by
TYPE:
|
control_guidance_start |
The percentage of total steps at which the ControlNet starts applying.
TYPE:
|
control_guidance_end |
The percentage of total steps at which the ControlNet stops applying.
TYPE:
|
original_size |
If
TYPE:
|
crops_coords_top_left |
TYPE:
|
target_size |
For most cases,
TYPE:
|
negative_original_size |
To negatively condition the generation process based on a specific image resolution. Part of SDXL's micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. For more information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208.
TYPE:
|
negative_crops_coords_top_left |
To negatively condition the generation process based on a specific crop coordinates. Part of SDXL's micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. For more information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208.
TYPE:
|
negative_target_size |
To negatively condition the generation process based on a target image resolution. It should be as same
as the
TYPE:
|
clip_skip |
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings.
TYPE:
|
callback_on_step_end |
A function or a subclass of
TYPE:
|
callback_on_step_end_tensor_inputs |
The list of tensor inputs for the
TYPE:
|
RETURNS | DESCRIPTION |
---|---|
[ |
Source code in mindone/diffusers/pipelines/controlnet_xs/pipeline_controlnet_xs_sd_xl.py
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|
mindone.diffusers.StableDiffusionXLControlNetXSPipeline.encode_prompt(prompt, prompt_2=None, num_images_per_prompt=1, do_classifier_free_guidance=True, negative_prompt=None, negative_prompt_2=None, prompt_embeds=None, negative_prompt_embeds=None, pooled_prompt_embeds=None, negative_pooled_prompt_embeds=None, lora_scale=None, clip_skip=None)
¶
Encodes the prompt into text encoder hidden states.
PARAMETER | DESCRIPTION |
---|---|
prompt |
prompt to be encoded
TYPE:
|
prompt_2 |
The prompt or prompts to be sent to the
TYPE:
|
num_images_per_prompt |
number of images that should be generated per prompt
TYPE:
|
do_classifier_free_guidance |
whether to use classifier free guidance or not
TYPE:
|
negative_prompt |
The prompt or prompts not to guide the image generation. If not defined, one has to pass
TYPE:
|
negative_prompt_2 |
The prompt or prompts not to guide the image generation to be sent to
TYPE:
|
prompt_embeds |
Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not
provided, text embeddings will be generated from
TYPE:
|
negative_prompt_embeds |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, negative_prompt_embeds will be generated from
TYPE:
|
pooled_prompt_embeds |
Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting.
If not provided, pooled text embeddings will be generated from
TYPE:
|
negative_pooled_prompt_embeds |
Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, pooled negative_prompt_embeds will be generated from
TYPE:
|
lora_scale |
A lora scale that will be applied to all LoRA layers of the text encoder if LoRA layers are loaded.
TYPE:
|
clip_skip |
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings.
TYPE:
|
Source code in mindone/diffusers/pipelines/controlnet_xs/pipeline_controlnet_xs_sd_xl.py
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|
mindone.diffusers.pipelines.stable_diffusion.StableDiffusionPipelineOutput
dataclass
¶
Bases: BaseOutput
Output class for Stable Diffusion pipelines.
Source code in mindone/diffusers/pipelines/stable_diffusion/pipeline_output.py
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