ControlNet with Stable Diffusion XL¶
ControlNet was introduced in Adding Conditional Control to Text-to-Image Diffusion Models by Lvmin Zhang, Anyi Rao, and Maneesh Agrawala.
With a ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. For example, if you provide a depth map, the ControlNet model generates an image that'll preserve the spatial information from the depth map. It is a more flexible and accurate way to control the image generation process.
The abstract from the paper is:
We present ControlNet, a neural network architecture to add spatial conditioning controls to large, pretrained text-to-image diffusion models. ControlNet locks the production-ready large diffusion models, and reuses their deep and robust encoding layers pretrained with billions of images as a strong backbone to learn a diverse set of conditional controls. The neural architecture is connected with "zero convolutions" (zero-initialized convolution layers) that progressively grow the parameters from zero and ensure that no harmful noise could affect the finetuning. We test various conditioning controls, eg, edges, depth, segmentation, human pose, etc, with Stable Diffusion, using single or multiple conditions, with or without prompts. We show that the training of ControlNets is robust with small (<50k) and large (>1m) datasets. Extensive results show that ControlNet may facilitate wider applications to control image diffusion models.
You can find additional smaller Stable Diffusion XL (SDXL) ControlNet checkpoints from the 🤗 Diffusers Hub organization, and browse community-trained checkpoints on the Hub.
Warning
🧪 Many of the SDXL ControlNet checkpoints are experimental, and there is a lot of room for improvement. Feel free to open an Issue and leave us feedback on how we can improve!
Tip
If you don't see a checkpoint you're interested in, you can train your own SDXL ControlNet with our training script.
Tip
Make sure to check out the Schedulers guide to learn how to explore the tradeoff between scheduler speed and quality, and see the reuse components across pipelines section to learn how to efficiently load the same components into multiple pipelines.
mindone.diffusers.StableDiffusionXLControlNetPipeline
¶
Bases: DiffusionPipeline
, StableDiffusionMixin
, TextualInversionLoaderMixin
, StableDiffusionXLLoraLoaderMixin
, IPAdapterMixin
, FromSingleFileMixin
Pipeline for text-to-image generation using Stable Diffusion XL with ControlNet guidance.
This model inherits from [DiffusionPipeline
]. Check the superclass documentation for the generic methods
implemented for all pipelines (downloading, saving, running on a particular device, etc.).
The pipeline also inherits the following loading methods
- [
~loaders.TextualInversionLoaderMixin.load_textual_inversion
] for loading textual inversion embeddings - [
~loaders.StableDiffusionXLLoraLoaderMixin.load_lora_weights
] for loading LoRA weights - [
~loaders.StableDiffusionXLLoraLoaderMixin.save_lora_weights
] for saving LoRA weights - [
~loaders.FromSingleFileMixin.from_single_file
] for loading.ckpt
files - [
~loaders.IPAdapterMixin.load_ip_adapter
] for loading IP Adapters
PARAMETER | DESCRIPTION |
---|---|
vae |
Variational Auto-Encoder (VAE) model to encode and decode images to and from latent representations.
TYPE:
|
text_encoder |
Frozen text-encoder (clip-vit-large-patch14).
TYPE:
|
text_encoder_2 |
Second frozen text-encoder (laion/CLIP-ViT-bigG-14-laion2B-39B-b160k).
TYPE:
|
tokenizer |
A
TYPE:
|
tokenizer_2 |
A
TYPE:
|
unet |
A
TYPE:
|
controlnet |
Provides additional conditioning to the
TYPE:
|
scheduler |
A scheduler to be used in combination with
TYPE:
|
force_zeros_for_empty_prompt |
Whether the negative prompt embeddings should always be set to 0. Also see the config of
TYPE:
|
add_watermarker |
Whether to use the invisible_watermark library to
watermark output images. If not defined, it defaults to
TYPE:
|
Source code in mindone/diffusers/pipelines/controlnet/pipeline_controlnet_sd_xl.py
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|
mindone.diffusers.StableDiffusionXLControlNetPipeline.__call__(prompt=None, prompt_2=None, image=None, height=None, width=None, num_inference_steps=50, timesteps=None, sigmas=None, denoising_end=None, guidance_scale=5.0, negative_prompt=None, negative_prompt_2=None, num_images_per_prompt=1, eta=0.0, generator=None, latents=None, prompt_embeds=None, negative_prompt_embeds=None, pooled_prompt_embeds=None, negative_pooled_prompt_embeds=None, ip_adapter_image=None, ip_adapter_image_embeds=None, output_type='pil', return_dict=False, cross_attention_kwargs=None, controlnet_conditioning_scale=1.0, guess_mode=False, control_guidance_start=0.0, control_guidance_end=1.0, original_size=None, crops_coords_top_left=(0, 0), target_size=None, negative_original_size=None, negative_crops_coords_top_left=(0, 0), negative_target_size=None, clip_skip=None, callback_on_step_end=None, callback_on_step_end_tensor_inputs=['latents'], **kwargs)
¶
The call function to the pipeline for generation.
PARAMETER | DESCRIPTION |
---|---|
prompt |
The prompt or prompts to guide image generation. If not defined, you need to pass
TYPE:
|
prompt_2 |
The prompt or prompts to be sent to
TYPE:
|
image |
The ControlNet input condition to provide guidance to the
TYPE:
|
height |
The height in pixels of the generated image. Anything below 512 pixels won't work well for stabilityai/stable-diffusion-xl-base-1.0 and checkpoints that are not specifically fine-tuned on low resolutions.
TYPE:
|
width |
The width in pixels of the generated image. Anything below 512 pixels won't work well for stabilityai/stable-diffusion-xl-base-1.0 and checkpoints that are not specifically fine-tuned on low resolutions.
TYPE:
|
num_inference_steps |
The number of denoising steps. More denoising steps usually lead to a higher quality image at the expense of slower inference.
TYPE:
|
timesteps |
Custom timesteps to use for the denoising process with schedulers which support a
TYPE:
|
sigmas |
Custom sigmas to use for the denoising process with schedulers which support a
TYPE:
|
denoising_end |
When specified, determines the fraction (between 0.0 and 1.0) of the total denoising process to be completed before it is intentionally prematurely terminated. As a result, the returned sample will still retain a substantial amount of noise as determined by the discrete timesteps selected by the scheduler. The denoising_end parameter should ideally be utilized when this pipeline forms a part of a "Mixture of Denoisers" multi-pipeline setup, as elaborated in Refining the Image Output
TYPE:
|
guidance_scale |
A higher guidance scale value encourages the model to generate images closely linked to the text
TYPE:
|
negative_prompt |
The prompt or prompts to guide what to not include in image generation. If not defined, you need to
pass
TYPE:
|
negative_prompt_2 |
The prompt or prompts to guide what to not include in image generation. This is sent to
TYPE:
|
num_images_per_prompt |
The number of images to generate per prompt.
TYPE:
|
eta |
Corresponds to parameter eta (η) from the DDIM paper. Only applies
to the [
TYPE:
|
generator |
A
TYPE:
|
latents |
Pre-generated noisy latents sampled from a Gaussian distribution, to be used as inputs for image
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents
tensor is generated by sampling using the supplied random
TYPE:
|
prompt_embeds |
Pre-generated text embeddings. Can be used to easily tweak text inputs (prompt weighting). If not
provided, text embeddings are generated from the
TYPE:
|
negative_prompt_embeds |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs (prompt weighting). If
not provided,
TYPE:
|
pooled_prompt_embeds |
Pre-generated pooled text embeddings. Can be used to easily tweak text inputs (prompt weighting). If
not provided, pooled text embeddings are generated from
TYPE:
|
negative_pooled_prompt_embeds |
Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs (prompt
weighting). If not provided, pooled
TYPE:
|
ip_adapter_image |
(
TYPE:
|
ip_adapter_image_embeds |
Pre-generated image embeddings for IP-Adapter. It should be a list of length same as number of
IP-adapters. Each element should be a tensor of shape
TYPE:
|
output_type |
The output format of the generated image. Choose between
TYPE:
|
return_dict |
Whether or not to return a [
TYPE:
|
cross_attention_kwargs |
A kwargs dictionary that if specified is passed along to the [
TYPE:
|
controlnet_conditioning_scale |
The outputs of the ControlNet are multiplied by
TYPE:
|
guess_mode |
The ControlNet encoder tries to recognize the content of the input image even if you remove all
prompts. A
TYPE:
|
control_guidance_start |
The percentage of total steps at which the ControlNet starts applying.
TYPE:
|
control_guidance_end |
The percentage of total steps at which the ControlNet stops applying.
TYPE:
|
original_size |
If
TYPE:
|
crops_coords_top_left |
TYPE:
|
target_size |
For most cases,
TYPE:
|
negative_original_size |
To negatively condition the generation process based on a specific image resolution. Part of SDXL's micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. For more information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208.
TYPE:
|
negative_crops_coords_top_left |
To negatively condition the generation process based on a specific crop coordinates. Part of SDXL's micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. For more information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208.
TYPE:
|
negative_target_size |
To negatively condition the generation process based on a target image resolution. It should be as same
as the
TYPE:
|
clip_skip |
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings.
TYPE:
|
callback_on_step_end |
A function or a subclass of
TYPE:
|
callback_on_step_end_tensor_inputs |
The list of tensor inputs for the
TYPE:
|
RETURNS | DESCRIPTION |
---|---|
[ |
Source code in mindone/diffusers/pipelines/controlnet/pipeline_controlnet_sd_xl.py
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|
mindone.diffusers.StableDiffusionXLControlNetPipeline.encode_prompt(prompt, prompt_2=None, num_images_per_prompt=1, do_classifier_free_guidance=True, negative_prompt=None, negative_prompt_2=None, prompt_embeds=None, negative_prompt_embeds=None, pooled_prompt_embeds=None, negative_pooled_prompt_embeds=None, lora_scale=None, clip_skip=None)
¶
Encodes the prompt into text encoder hidden states.
PARAMETER | DESCRIPTION |
---|---|
prompt |
prompt to be encoded
TYPE:
|
prompt_2 |
The prompt or prompts to be sent to the
TYPE:
|
num_images_per_prompt |
number of images that should be generated per prompt
TYPE:
|
do_classifier_free_guidance |
whether to use classifier free guidance or not
TYPE:
|
negative_prompt |
The prompt or prompts not to guide the image generation. If not defined, one has to pass
TYPE:
|
negative_prompt_2 |
The prompt or prompts not to guide the image generation to be sent to
TYPE:
|
prompt_embeds |
Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not
provided, text embeddings will be generated from
TYPE:
|
negative_prompt_embeds |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, negative_prompt_embeds will be generated from
TYPE:
|
pooled_prompt_embeds |
Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting.
If not provided, pooled text embeddings will be generated from
TYPE:
|
negative_pooled_prompt_embeds |
Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, pooled negative_prompt_embeds will be generated from
TYPE:
|
lora_scale |
A lora scale that will be applied to all LoRA layers of the text encoder if LoRA layers are loaded.
TYPE:
|
clip_skip |
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings.
TYPE:
|
Source code in mindone/diffusers/pipelines/controlnet/pipeline_controlnet_sd_xl.py
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|
mindone.diffusers.StableDiffusionXLControlNetPipeline.get_guidance_scale_embedding(w, embedding_dim=512, dtype=ms.float32)
¶
See https://github.com/google-research/vdm/blob/dc27b98a554f65cdc654b800da5aa1846545d41b/model_vdm.py#L298
PARAMETER | DESCRIPTION |
---|---|
w |
Generate embedding vectors with a specified guidance scale to subsequently enrich timestep embeddings.
TYPE:
|
embedding_dim |
Dimension of the embeddings to generate.
TYPE:
|
dtype |
Data type of the generated embeddings.
TYPE:
|
RETURNS | DESCRIPTION |
---|---|
Tensor
|
|
Source code in mindone/diffusers/pipelines/controlnet/pipeline_controlnet_sd_xl.py
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|
mindone.diffusers.StableDiffusionXLControlNetImg2ImgPipeline
¶
Bases: DiffusionPipeline
, StableDiffusionMixin
, TextualInversionLoaderMixin
, StableDiffusionXLLoraLoaderMixin
, FromSingleFileMixin
, IPAdapterMixin
Pipeline for image-to-image generation using Stable Diffusion XL with ControlNet guidance.
This model inherits from [DiffusionPipeline
]. Check the superclass documentation for the generic methods the
library implements for all the pipelines (such as downloading or saving, running on a particular device, etc.)
The pipeline also inherits the following loading methods
- [
~loaders.TextualInversionLoaderMixin.load_textual_inversion
] for loading textual inversion embeddings - [
~loaders.StableDiffusionXLLoraLoaderMixin.load_lora_weights
] for loading LoRA weights - [
~loaders.StableDiffusionXLLoraLoaderMixin.save_lora_weights
] for saving LoRA weights - [
~loaders.IPAdapterMixin.load_ip_adapter
] for loading IP Adapters
PARAMETER | DESCRIPTION |
---|---|
vae |
Variational Auto-Encoder (VAE) Model to encode and decode images to and from latent representations.
TYPE:
|
text_encoder |
Frozen text-encoder. Stable Diffusion uses the text portion of CLIP, specifically the clip-vit-large-patch14 variant.
TYPE:
|
text_encoder_2 |
Second frozen text-encoder. Stable Diffusion XL uses the text and pool portion of CLIP, specifically the laion/CLIP-ViT-bigG-14-laion2B-39B-b160k variant.
TYPE:
|
tokenizer |
Tokenizer of class CLIPTokenizer.
TYPE:
|
tokenizer_2 |
Second Tokenizer of class CLIPTokenizer.
TYPE:
|
unet |
Conditional U-Net architecture to denoise the encoded image latents.
TYPE:
|
controlnet |
Provides additional conditioning to the unet during the denoising process. If you set multiple ControlNets as a list, the outputs from each ControlNet are added together to create one combined additional conditioning.
TYPE:
|
scheduler |
A scheduler to be used in combination with
TYPE:
|
requires_aesthetics_score |
Whether the
TYPE:
|
force_zeros_for_empty_prompt |
Whether the negative prompt embeddings shall be forced to always be set to 0. Also see the config of
TYPE:
|
add_watermarker |
Whether to use the invisible_watermark library to watermark output images. If not defined, it will default to True if the package is installed, otherwise no watermarker will be used.
TYPE:
|
feature_extractor |
A
TYPE:
|
Source code in mindone/diffusers/pipelines/controlnet/pipeline_controlnet_sd_xl_img2img.py
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|
mindone.diffusers.StableDiffusionXLControlNetImg2ImgPipeline.__call__(prompt=None, prompt_2=None, image=None, control_image=None, height=None, width=None, strength=0.8, num_inference_steps=50, guidance_scale=5.0, negative_prompt=None, negative_prompt_2=None, num_images_per_prompt=1, eta=0.0, generator=None, latents=None, prompt_embeds=None, negative_prompt_embeds=None, pooled_prompt_embeds=None, negative_pooled_prompt_embeds=None, ip_adapter_image=None, ip_adapter_image_embeds=None, output_type='pil', return_dict=False, cross_attention_kwargs=None, controlnet_conditioning_scale=0.8, guess_mode=False, control_guidance_start=0.0, control_guidance_end=1.0, original_size=None, crops_coords_top_left=(0, 0), target_size=None, negative_original_size=None, negative_crops_coords_top_left=(0, 0), negative_target_size=None, aesthetic_score=6.0, negative_aesthetic_score=2.5, clip_skip=None, callback_on_step_end=None, callback_on_step_end_tensor_inputs=['latents'], **kwargs)
¶
Function invoked when calling the pipeline for generation.
PARAMETER | DESCRIPTION |
---|---|
prompt |
The prompt or prompts to guide the image generation. If not defined, one has to pass
TYPE:
|
prompt_2 |
The prompt or prompts to be sent to the
TYPE:
|
image |
The initial image will be used as the starting point for the image generation process. Can also accept
image latents as
TYPE:
|
control_image |
The ControlNet input condition. ControlNet uses this input condition to generate guidance to Unet. If
the type is specified as
TYPE:
|
height |
The height in pixels of the generated image. Anything below 512 pixels won't work well for stabilityai/stable-diffusion-xl-base-1.0 and checkpoints that are not specifically fine-tuned on low resolutions.
TYPE:
|
width |
The width in pixels of the generated image. Anything below 512 pixels won't work well for stabilityai/stable-diffusion-xl-base-1.0 and checkpoints that are not specifically fine-tuned on low resolutions.
TYPE:
|
strength |
Indicates extent to transform the reference
TYPE:
|
num_inference_steps |
The number of denoising steps. More denoising steps usually lead to a higher quality image at the expense of slower inference.
TYPE:
|
guidance_scale |
Guidance scale as defined in Classifier-Free Diffusion Guidance.
TYPE:
|
negative_prompt |
The prompt or prompts not to guide the image generation. If not defined, one has to pass
TYPE:
|
negative_prompt_2 |
The prompt or prompts not to guide the image generation to be sent to
TYPE:
|
num_images_per_prompt |
The number of images to generate per prompt.
TYPE:
|
eta |
Corresponds to parameter eta (η) in the DDIM paper: https://arxiv.org/abs/2010.02502. Only applies to
[
TYPE:
|
generator |
One or a list of torch generator(s) to make generation deterministic.
TYPE:
|
latents |
Pre-generated noisy latents, sampled from a Gaussian distribution, to be used as inputs for image
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents
tensor will ge generated by sampling using the supplied random
TYPE:
|
prompt_embeds |
Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not
provided, text embeddings will be generated from
TYPE:
|
negative_prompt_embeds |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, negative_prompt_embeds will be generated from
TYPE:
|
pooled_prompt_embeds |
Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting.
If not provided, pooled text embeddings will be generated from
TYPE:
|
negative_pooled_prompt_embeds |
Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, pooled negative_prompt_embeds will be generated from
TYPE:
|
ip_adapter_image |
(
TYPE:
|
ip_adapter_image_embeds |
Pre-generated image embeddings for IP-Adapter. It should be a list of length same as number of
IP-adapters. Each element should be a tensor of shape
TYPE:
|
output_type |
The output format of the generate image. Choose between
PIL:
TYPE:
|
return_dict |
Whether or not to return a [
TYPE:
|
cross_attention_kwargs |
A kwargs dictionary that if specified is passed along to the
TYPE:
|
controlnet_conditioning_scale |
The outputs of the controlnet are multiplied by
TYPE:
|
guess_mode |
In this mode, the ControlNet encoder will try best to recognize the content of the input image even if
you remove all prompts. The
TYPE:
|
control_guidance_start |
The percentage of total steps at which the controlnet starts applying.
TYPE:
|
control_guidance_end |
The percentage of total steps at which the controlnet stops applying.
TYPE:
|
original_size |
If
TYPE:
|
crops_coords_top_left |
TYPE:
|
target_size |
For most cases,
TYPE:
|
negative_original_size |
To negatively condition the generation process based on a specific image resolution. Part of SDXL's micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. For more information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208.
TYPE:
|
negative_crops_coords_top_left |
To negatively condition the generation process based on a specific crop coordinates. Part of SDXL's micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. For more information, refer to this issue thread: https://github.com/huggingface/diffusers/issues/4208.
TYPE:
|
negative_target_size |
To negatively condition the generation process based on a target image resolution. It should be as same
as the
TYPE:
|
aesthetic_score |
Used to simulate an aesthetic score of the generated image by influencing the positive text condition. Part of SDXL's micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952.
TYPE:
|
negative_aesthetic_score |
Part of SDXL's micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. Can be used to simulate an aesthetic score of the generated image by influencing the negative text condition.
TYPE:
|
clip_skip |
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings.
TYPE:
|
callback_on_step_end |
A function or a subclass of
TYPE:
|
callback_on_step_end_tensor_inputs |
The list of tensor inputs for the
TYPE:
|
RETURNS | DESCRIPTION |
---|---|
[ |
|
[ |
|
containing the output images. |
Source code in mindone/diffusers/pipelines/controlnet/pipeline_controlnet_sd_xl_img2img.py
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|
mindone.diffusers.StableDiffusionXLControlNetImg2ImgPipeline.encode_prompt(prompt, prompt_2=None, num_images_per_prompt=1, do_classifier_free_guidance=True, negative_prompt=None, negative_prompt_2=None, prompt_embeds=None, negative_prompt_embeds=None, pooled_prompt_embeds=None, negative_pooled_prompt_embeds=None, lora_scale=None, clip_skip=None)
¶
Encodes the prompt into text encoder hidden states.
PARAMETER | DESCRIPTION |
---|---|
prompt |
prompt to be encoded
TYPE:
|
prompt_2 |
The prompt or prompts to be sent to the
TYPE:
|
num_images_per_prompt |
number of images that should be generated per prompt
TYPE:
|
do_classifier_free_guidance |
whether to use classifier free guidance or not
TYPE:
|
negative_prompt |
The prompt or prompts not to guide the image generation. If not defined, one has to pass
TYPE:
|
negative_prompt_2 |
The prompt or prompts not to guide the image generation to be sent to
TYPE:
|
prompt_embeds |
Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not
provided, text embeddings will be generated from
TYPE:
|
negative_prompt_embeds |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, negative_prompt_embeds will be generated from
TYPE:
|
pooled_prompt_embeds |
Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting.
If not provided, pooled text embeddings will be generated from
TYPE:
|
negative_pooled_prompt_embeds |
Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, pooled negative_prompt_embeds will be generated from
TYPE:
|
lora_scale |
A lora scale that will be applied to all LoRA layers of the text encoder if LoRA layers are loaded.
TYPE:
|
clip_skip |
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings.
TYPE:
|
Source code in mindone/diffusers/pipelines/controlnet/pipeline_controlnet_sd_xl_img2img.py
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|
mindone.diffusers.StableDiffusionXLControlNetInpaintPipeline
¶
Bases: DiffusionPipeline
, StableDiffusionMixin
, StableDiffusionXLLoraLoaderMixin
, FromSingleFileMixin
, IPAdapterMixin
, TextualInversionLoaderMixin
Pipeline for text-to-image generation using Stable Diffusion XL.
This model inherits from [DiffusionPipeline
]. Check the superclass documentation for the generic methods the
library implements for all the pipelines (such as downloading or saving, running on a particular device, etc.)
The pipeline also inherits the following loading methods
- [
~loaders.TextualInversionLoaderMixin.load_textual_inversion
] for loading textual inversion embeddings - [
~loaders.StableDiffusionXLLoraLoaderMixin.load_lora_weights
] for loading LoRA weights - [
~loaders.StableDiffusionXLLoraLoaderMixin.save_lora_weights
] for saving LoRA weights - [
~loaders.FromSingleFileMixin.from_single_file
] for loading.ckpt
files - [
~loaders.IPAdapterMixin.load_ip_adapter
] for loading IP Adapters
PARAMETER | DESCRIPTION |
---|---|
vae |
Variational Auto-Encoder (VAE) Model to encode and decode images to and from latent representations.
TYPE:
|
text_encoder |
Frozen text-encoder. Stable Diffusion XL uses the text portion of CLIP, specifically the clip-vit-large-patch14 variant.
TYPE:
|
text_encoder_2 |
Second frozen text-encoder. Stable Diffusion XL uses the text and pool portion of CLIP, specifically the laion/CLIP-ViT-bigG-14-laion2B-39B-b160k variant.
TYPE:
|
tokenizer |
Tokenizer of class CLIPTokenizer.
TYPE:
|
tokenizer_2 |
Second Tokenizer of class CLIPTokenizer.
TYPE:
|
unet |
Conditional U-Net architecture to denoise the encoded image latents.
TYPE:
|
scheduler |
A scheduler to be used in combination with
TYPE:
|
Source code in mindone/diffusers/pipelines/controlnet/pipeline_controlnet_inpaint_sd_xl.py
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mindone.diffusers.StableDiffusionXLControlNetInpaintPipeline.__call__(prompt=None, prompt_2=None, image=None, mask_image=None, control_image=None, height=None, width=None, padding_mask_crop=None, strength=0.9999, num_inference_steps=50, denoising_start=None, denoising_end=None, guidance_scale=5.0, negative_prompt=None, negative_prompt_2=None, num_images_per_prompt=1, eta=0.0, generator=None, latents=None, prompt_embeds=None, negative_prompt_embeds=None, ip_adapter_image=None, ip_adapter_image_embeds=None, pooled_prompt_embeds=None, negative_pooled_prompt_embeds=None, output_type='pil', return_dict=False, cross_attention_kwargs=None, controlnet_conditioning_scale=1.0, guess_mode=False, control_guidance_start=0.0, control_guidance_end=1.0, guidance_rescale=0.0, original_size=None, crops_coords_top_left=(0, 0), target_size=None, aesthetic_score=6.0, negative_aesthetic_score=2.5, clip_skip=None, callback_on_step_end=None, callback_on_step_end_tensor_inputs=['latents'], **kwargs)
¶
Function invoked when calling the pipeline for generation.
PARAMETER | DESCRIPTION |
---|---|
prompt |
The prompt or prompts to guide the image generation. If not defined, one has to pass
TYPE:
|
prompt_2 |
The prompt or prompts to be sent to the
TYPE:
|
image |
TYPE:
|
mask_image |
TYPE:
|
height |
The height in pixels of the generated image.
TYPE:
|
width |
The width in pixels of the generated image.
TYPE:
|
padding_mask_crop |
The size of margin in the crop to be applied to the image and masking. If
TYPE:
|
strength |
Conceptually, indicates how much to transform the masked portion of the reference
TYPE:
|
num_inference_steps |
The number of denoising steps. More denoising steps usually lead to a higher quality image at the expense of slower inference.
TYPE:
|
denoising_start |
When specified, indicates the fraction (between 0.0 and 1.0) of the total denoising process to be
bypassed before it is initiated. Consequently, the initial part of the denoising process is skipped and
it is assumed that the passed
TYPE:
|
denoising_end |
When specified, determines the fraction (between 0.0 and 1.0) of the total denoising process to be
completed before it is intentionally prematurely terminated. As a result, the returned sample will
still retain a substantial amount of noise (ca. final 20% of timesteps still needed) and should be
denoised by a successor pipeline that has
TYPE:
|
guidance_scale |
Guidance scale as defined in Classifier-Free Diffusion Guidance.
TYPE:
|
negative_prompt |
The prompt or prompts not to guide the image generation. If not defined, one has to pass
TYPE:
|
negative_prompt_2 |
The prompt or prompts not to guide the image generation to be sent to
TYPE:
|
prompt_embeds |
Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not
provided, text embeddings will be generated from
TYPE:
|
negative_prompt_embeds |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, negative_prompt_embeds will be generated from
TYPE:
|
ip_adapter_image |
(
TYPE:
|
ip_adapter_image_embeds |
Pre-generated image embeddings for IP-Adapter. It should be a list of length same as number of
IP-adapters. Each element should be a tensor of shape
TYPE:
|
pooled_prompt_embeds |
Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting.
If not provided, pooled text embeddings will be generated from
TYPE:
|
negative_pooled_prompt_embeds |
Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, pooled negative_prompt_embeds will be generated from
TYPE:
|
num_images_per_prompt |
The number of images to generate per prompt.
TYPE:
|
eta |
Corresponds to parameter eta (η) in the DDIM paper: https://arxiv.org/abs/2010.02502. Only applies to
[
TYPE:
|
generator |
One or a list of torch generator(s) to make generation deterministic.
TYPE:
|
latents |
Pre-generated noisy latents, sampled from a Gaussian distribution, to be used as inputs for image
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents
tensor will ge generated by sampling using the supplied random
TYPE:
|
output_type |
The output format of the generate image. Choose between
PIL:
TYPE:
|
return_dict |
Whether or not to return a [
TYPE:
|
cross_attention_kwargs |
A kwargs dictionary that if specified is passed along to the
TYPE:
|
original_size |
If
TYPE:
|
crops_coords_top_left |
TYPE:
|
target_size |
For most cases,
TYPE:
|
aesthetic_score |
Used to simulate an aesthetic score of the generated image by influencing the positive text condition. Part of SDXL's micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952.
TYPE:
|
negative_aesthetic_score |
Part of SDXL's micro-conditioning as explained in section 2.2 of https://huggingface.co/papers/2307.01952. Can be used to simulate an aesthetic score of the generated image by influencing the negative text condition.
TYPE:
|
clip_skip |
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings.
TYPE:
|
callback_on_step_end |
A function or a subclass of
TYPE:
|
callback_on_step_end_tensor_inputs |
The list of tensor inputs for the
TYPE:
|
RETURNS | DESCRIPTION |
---|---|
[ |
|
[ |
|
|
Source code in mindone/diffusers/pipelines/controlnet/pipeline_controlnet_inpaint_sd_xl.py
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mindone.diffusers.StableDiffusionXLControlNetInpaintPipeline.encode_prompt(prompt, prompt_2=None, num_images_per_prompt=1, do_classifier_free_guidance=True, negative_prompt=None, negative_prompt_2=None, prompt_embeds=None, negative_prompt_embeds=None, pooled_prompt_embeds=None, negative_pooled_prompt_embeds=None, lora_scale=None, clip_skip=None)
¶
Encodes the prompt into text encoder hidden states.
PARAMETER | DESCRIPTION |
---|---|
prompt |
prompt to be encoded
TYPE:
|
prompt_2 |
The prompt or prompts to be sent to the
TYPE:
|
num_images_per_prompt |
number of images that should be generated per prompt
TYPE:
|
do_classifier_free_guidance |
whether to use classifier free guidance or not
TYPE:
|
negative_prompt |
The prompt or prompts not to guide the image generation. If not defined, one has to pass
TYPE:
|
negative_prompt_2 |
The prompt or prompts not to guide the image generation to be sent to
TYPE:
|
prompt_embeds |
Pre-generated text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting. If not
provided, text embeddings will be generated from
TYPE:
|
negative_prompt_embeds |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, negative_prompt_embeds will be generated from
TYPE:
|
pooled_prompt_embeds |
Pre-generated pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt weighting.
If not provided, pooled text embeddings will be generated from
TYPE:
|
negative_pooled_prompt_embeds |
Pre-generated negative pooled text embeddings. Can be used to easily tweak text inputs, e.g. prompt
weighting. If not provided, pooled negative_prompt_embeds will be generated from
TYPE:
|
lora_scale |
A lora scale that will be applied to all LoRA layers of the text encoder if LoRA layers are loaded.
TYPE:
|
clip_skip |
Number of layers to be skipped from CLIP while computing the prompt embeddings. A value of 1 means that the output of the pre-final layer will be used for computing the prompt embeddings.
TYPE:
|
Source code in mindone/diffusers/pipelines/controlnet/pipeline_controlnet_inpaint_sd_xl.py
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mindone.diffusers.pipelines.stable_diffusion.StableDiffusionPipelineOutput
dataclass
¶
Bases: BaseOutput
Output class for Stable Diffusion pipelines.
Source code in mindone/diffusers/pipelines/stable_diffusion/pipeline_output.py
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